List II
List II
List of Courses
Algebraic Geometry
Algebraic Topology
Analysis of Functions
Applications of Quantum Mechanics
Applied Probability
Asymptotic Methods
Automata and Formal Languages
Classical Dynamics
Coding and Cryptography
Cosmology
Differential Geometry
Dynamical Systems
Electrodynamics
Fluid Dynamics
Further Complex Methods
Galois Theory
General Relativity
Graph Theory
Integrable Systems
Linear Analysis
Logic and Set Theory
Mathematical Biology
Mathematics of Machine Learning
Number Fields
Number Theory
Numerical Analysis
Principles of Quantum Mechanics
Principles of Statistics
Probability and Measure
Quantum Information and Computation
Representation Theory
2
Riemann Surfaces
Statistical Modelling
Statistical Physics
Stochastic Financial Models
Topics in Analysis
Waves
Z := {(x, ex ) | x ∈ C} ⊆ A2
Paper 2, Section II
25F Algebraic Geometry
In this question, all algebraic varieties are over an algebraically closed field k.
Let X be an affine variety. Define the tangent space of X at a point P ∈ X. Define
the dimension of X in terms of (i) the tangent spaces of X and (ii) Krull dimension. Say
what it means for the variety to be singular at P .
Assume the characteristic of the field k is not 2. Let X := Z(x21 − x32 , x23 − x34 ) ⊆ A4 .
Calculate the tangent space of X at each point of X.
Consider the subset Y ⊆ P4 consisting of points with homogeneous coordinates
(y0 : y1 : y2 : y3 : y4 ) such that the matrix
y0 y1 y2
y2 y3 y4
has rank one. Show that Y is a closed subset of P4 in the Zariski topology. You may now
assume Y is irreducible in the Zariski topology, and hence is a projective variety. What is
the dimension of Y ? Show that Y is non-singular.
X = Z(y 2 z − x3 + 4xz 2 − z 3 ) ⊆ P2
Paper 4, Section II
24F Algebraic Geometry
In this question, all algebraic varieties are over the complex numbers C.
Let φ : P2 → P1 be a morphism. Show that φ is constant. [You may use without
proof the fact that any two curves in P2 have non-empty intersection.]
Show that there is no closed subvariety of P1 × P1 × P1 isomorphic to P2 .
State the Riemann–Hurwitz theorem.
Consider a non-singular projective curve X ⊆ P1 × P1 defined by a bihomogeneous
equation f (x0 , x1 , y0 , y1 ) = 0 which is homogeneous of degree 2 in x0 , x1 and homogeneous
of degree 3 in y0 , y1 . Here x0 , x1 are coordinates on the first P1 and y0 , y1 are coordinates
on the second P1 . Compute the genus of X. Deduce that X is not isomorphic to a non-
singular projective curve in P2 . [You may use without proof the fact that a non-singular
projective curve in P2 of degree d has genus g = (d − 1)(d − 2)/2.]
(ii) p ◦ γ̃ = γ.
Y ×Z
p : Yd
/ϕ := −→ Y /ϕ
(a, i) ∼ (ϕ(a), i − 1) for a ∈ A, i ∈ Z
[(y, i)] 7−→ [y]
Paper 2, Section II
21J Algebraic Topology
State the Seifert–van Kampen theorem.
If (X, x0 ) is a based topological space, and f : (S n−1 , ∗) → (X, x0 ) is a map of based
spaces, define the space X ∪f Dn obtained by attaching an n-cell to X along f . For n = 2,
carefully prove a formula describing π1 (X ∪f D2 , x0 ) in terms of the group π1 (X, x0 ) and
the element [f ] ∈ π1 (X, x0 ).
Writing S 1 ∨ S 1 for the wedge of two circles, calculate π1 (S 1 ∨ S 1 , ∗).
Explain how to attach 2-cells to S 1 ∨ S 1 to obtain a space whose fundamental group
is the symmetric group on 3 letters, proving carefully that this is indeed the fundamental
group obtained.
[You may use any description of the group π1 (S 1 , ∗), provided it is clearly stated.
You should justify any presentation of the symmetric group on 3 letters that you use.]
Ck (L)
Ck (L, K) := ,
Ck (K)
show that the boundary map of L descends to give C• (L, K) the structure of a chain
complex. Describe a long exact sequence relating H∗ (K), H∗ (L), and the homology
H∗ (L, K) of C• (L, K).
In the following, ∆n+1 denotes the standard (n + 1)-simplex and ∂∆n+1 denotes its
boundary, i.e. the union of all its simplices of dimension < n + 1. Suppose that n ⩾ 2.
Using that ∆n+1 has the same homology as a point, calculate the homology of
∂∆n+1 .
Recall that the rank rk A of an abelian group A is the maximal r such that Zr ⩽ A.
If K ⩽ ∂∆n+1 is a sub-simplicial complex, by considering H∗ (∂∆n+1 , K) show that for
each 0 < k ⩽ n we have
n+2
rk Hk (K) ⩽ − #(k + 1)-simplices of K.
k+2
[You may use that the rank of a subgroup or quotient of A is at most that of A.]
If K < ∂∆n+1 is a proper sub-simplicial complex, show that Hn (K) = 0.
Paper 4, Section II
21J Algebraic Topology
State the Mayer–Vietoris theorem for a simplicial complex K which is the union of
two subcomplexes M and N .
For each k = 2, 3, 4, . . ., construct an explicit simplicial complex K having
Z
i=0
∼
Hi (K) = Z/k i = 1
0 else.
∥T f ∥Lp ⩽ ∥f ∥Lp .
[Hint: Consider fz (y) = f (yz) and show first that ∥fz ∥Lp = z −1/p ∥f ∥Lp .]
[You may use the identity
nZ o
∥f ∥Lp = sup |f (x)g(x)|dx : g ∈ Lq , ∥g∥Lq ⩽ 1 , q conjugate to p,
without proof.]
Let I = [0, 1] and denote by L∞ the Banach space of µ-essentially bounded functions
on I, where µ is Lebesgue measure. Show that (L∞ )′ does not coincide with L1 (µ). [Hint:
Extend the functional ℓ(f ) = f (0) from the subspace C(I) of continuous functions on I to
L∞ .]
Paper 3, Section II
22G Analysis of Functions
State and prove the Sobolev embedding theorem H s (Rn ) ⊂ L∞ (Rn ).
Show that H 1 (R3 ) contains an unbounded function.
[You may use results from Fourier analysis without proof, provided they are carefully
stated.]
Paper 4, Section II
23G Analysis of Functions
(a) For f ∈ H s (Rn ), s ∈ R, show that there exists a unique weak solution u ∈ S ′ (Rn )
to the equation
∆2 u + u = f. (⋆)
(b) Show that if f has compact support in Rn and is infinitely differentiable, then
the solution u is also infinitely differentiable.
(c) Now let n = 3 and let f = δ0 be the Dirac measure at zero. Show that there
exists a unique continuous function u solving (⋆).
[You may use results about Fourier transforms and Sobolev spaces from the course
without proof, provided they are clearly stated.]
where Pl (cos θ) are Legendre polynomials, and boundary condition ul (0) = 0, the time-
independent Schrödinger equation for the wavefunction ψ(r) of the particle reduces to
d2 l(l + 1) 2m 2 ℏ2 k 2
− 2+ + V (r) ul (r) = k ul (r) , with E = .
dr r2 ℏ2 2m
(a) The asymptotic behaviour, for large r, of the wavefunction can be written
∞
X −ikr
2l + 1 l+1 e eikr
ψ(r) ∼ (−1) + Sl (k) Pl (cos θ) .
2ik r r
l=0
Hence deduce that Sl (k) = e2iδl (k) for some real function δl (k), and that δl (k) = −δl (−k).
(b) Focus now on the low-momentum behaviour, where the l = 0 mode dominates.
For some choice of potential V (r),
(k + 3iλ)(k + 2iλ)
S0 (k) = ,
(k − 3iλ)(k − 2iλ)
where λ is a real positive constant. Evaluate the scattering length as , and give an estimate
for the total cross-section σT at low energies.
Briefly explain the significance of the poles of S0 (k). Are there any resonances in
S0 (k)? Why is it important that λ is positive?
where αn are complex variational parameters, and the |ϕn ⟩ form an orthonormal set, i.e.,
⟨ϕm |ϕn ⟩ = δmn , for m and n = 1, 2, . . . N .
Apply the variational method to H with trial wavefunction |ψ⟩, and show that the
lowest eigenvalue of the matrix H, which has entries Hnm = ⟨ϕn |H|ϕm ⟩, gives the optimal
upper bound on the ground state energy E0 .
(c) Consider a particle of mass m in an infinite one-dimensional square well of width
a, with a linear potential in the well,
( x
V0 0 ⩽ x ⩽ a,
V (x) = a
∞ otherwise ,
9ℏ2
where V0 = . Determine an upper bound for the ground state energy of this system
ma2
using ( q
2
sin nπx 0 ⩽ x ⩽ a,
ϕn = a a
0 otherwise ,
with n = 1, 2, as trial wavefunctions.
ℏ2 d2 ψ
− + V (x)ψ(x) = E ψ(x) , (⋆)
2m dx2
where E is the energy of the system. The Floquet matrix F (E) is given by
ψ1 (x + a) ψ1 (x)
= F (E) ,
ψ2 (x + a) ψ2 (x)
(iii) Explain what the implications on the band structure are when (TrF )2 < 4
and when (TrF )2 > 4.
(b) Consider a tight-binding Hamiltonian that acts upon a single band of localised
states in one dimension,
X
Ht = t |j⟩⟨j + 1| + |j⟩⟨j − 1| + (−1)j |j⟩⟨j| ,
j∈Z
where t is a positive constant. The integer j should be thought of as indexing sites along a
chain of atoms separated from each other by a distance a; the state |j⟩ locates an electron
on atom j.
(i) Consider the translation operator Tℓ which acts on the states |j⟩ by shifting
them to another site, Tℓ |j⟩ = |j + ℓ⟩, with ℓ ∈ Z. Determine the values of
ℓ for which Tℓ commutes with the Hamiltonian Ht . Hence determine the
lattice spacing of the system.
(ii) Using Bloch’s theorem, show that the Hamiltonian of the system can be
reduced to a 2×2 matrix. Find the eigenvalues of this matrix.
(iii) What is the range of the first Brillouin zone? Plot the energy bands for
the first Brillouin zone. In your plot you should clearly label the axes and
indicate the values of the energies at the boundaries of the Brillouin zone.
What is the bandwidth of each of the bands?
where A is the vector potential and ϕ is the electric potential, with B = ∇ × A and
E = −∇ϕ − ∂A∂t .
(a) The Schrödinger equation (⋆) should be invariant under the gauge transforma-
tions
A(x, t) → A(x, t) + ∇f (x, t) ,
∂
ϕ(x, t) → ϕ(x, t) − f (x, t) ,
∂t
where f is an arbitrary smooth function, together with a suitable transformation of ψ(x, t).
State such a suitable transformation of ψ(x, t) and check explicitly that (⋆) is invariant
under these transformations.
(b) In the presence of electric and magnetic fields, the probability current is modified
to
iℏ e
J (x, t) = −(ψ ∗ ∇ψ − ψ∇ψ ∗ ) + ψ ∗ ψA .
2m m
Show that J is gauge invariant and that it satisfies the conservation equation
∂ρ
+∇·J =0 ,
∂t
where ψ(x, t) satisfies the Schrödinger equation and ρ = ψ ∗ ψ.
(c) Consider the situation with constant electromagnetic fields, where B = (0, 0, B),
E = (0, E, 0), and B and E are constants. Verify that the vector and electric potential can
be chosen to have the form
A = B(−y, 0, 0) , ϕ = −Ey .
Restrict the motion to the (x, y) plane. Show that stationary states are given by
where n is a non-negative integer, and you should determine both ω1 and W explicitly.
Briefly discuss what happens to the Landau levels in the presence of this electric field.
Part II, 2024 List of Questions
13
qin P(Jn ⩽ t < Jn+1 |Y0 = i0 , . . . , Yn = in ) = qi0 P(Jn ⩽ t < Jn+1 |Y0 = in , . . . , Yn = i0 ) .
(b) Now let (Xt )t⩾0 be irreducible. Fix any h > 0 and let Zn = Xnh for
n = 0, 1, 2, . . .. Show that (Zn )n⩾0 is a discrete-time Markov chain and give its transition
matrix. Show that (Xt )t⩾0 is recurrent if and only if (Zn )n⩾0 is recurrent.
(c) Finally, let I be finite and f : I → R be a function, identified with the vector
(f (x))x∈I . Show that
Ex (f (Xt )) − f (x)
Qf (x) = lim ,
t→0+ t
and Z t
Ex (f (Xt )) = f (x) + Ex (Qf (Xs )) ds.
0
(ii) Let n ⩾ 1. Show that conditional on the event {Xt = n + 1}, the times of
the n births have the distribution of the order statistics of n i.i.d. random
variables with probability density function
λeλx
f (x) = , 0 ⩽ x ⩽ t.
eλt − 1
(b) Let Xt be a birth and death process with rates λn = nλ and µn = nµ for n ∈ N,
λ > 0, µ > 0, and assume that X0 = 1. Let h(t) = P(Xt = 0).
Paper 3, Section II
27K Applied Probability
(a) State and prove Little’s formula for a regenerative process.
(b) What is an M/G/1 queue? Find the expected length of the busy period for an
M/G/1 queue.
(c) Suppose customers arrive at a single server at rate λ > 0 and require an
exponential amount of service time with rate µ > 0. Customers not being served are
impatient and will leave at rate δ > 0, independently of their position in the queue. Let
Xt denote the length of the queue at time t.
(i) Show that the system X = (Xt )t⩾0 has an invariant distribution.
show that
Pn−1
f (x) f (x) − k=0 ak ϕk (x)
a0 = lim and an = lim . (†)
x→x0 ϕ0 (x) x→x0 ϕn (x)
(b) Consider the asymptotic sequence ϕn (x), defined by ϕ0 = x−1 and ϕn (x) =
x−n+1 e−x for n ⩾ 1, as x → ∞, and the function
1 xe−x
f (x) = + .
x x−1
Find the asymptotic expansion of f (x) with respect to the ϕn (x) as x → ∞.
Verify explicitly that your coefficients satisfy (†) for all n.
What is the asymptotic expansion of f (x) with respect to the asymptotic sequence
ψn (x) = x−n as x → ∞?
(c) Consider the sine-integral function,
Z ∞
sin(xt)
si(x) = dt.
1 t
defined for x ∈ [1, ∞), with eigenvalue λ > 1 and small parameter 0 < ϵ ≪ 1, and with
boundary conditions y(1) = 0 and y → 0 as x → ∞.
for integer n. [You may quote the following asymptotic behaviour of the Airy function:
h i
Ai(t) ∼ π −1/2 (−t)−1/4 cos 23 (−t)3/2 − 14 π as t → −∞. ]
S S
A S A A
a b A A a a b b
a b b b a b
(ii) w = abbbaa.
(c) Prove that the language L := {an bm amin(n,m) : n, m > 1} is not context-free.
[You may use the context-free pumping lemma without proof.]
Paper 2, Section I
4J Automata and Formal Languages
(a) Define what an index set is and when it is called non-trivial.
(b) Define the index set Inf .
(c) State Rice’s theorem.
(d) Let X ⊆ W and let Inf X := {w ∈ Inf : ran(fw,1 ) ⊆ X}. Show, by modifying
the proof of Rice’s theorem or otherwise, that for each nonempty X, the set Inf X is not
computable.
Q := P ∪ {S → ε} and
H := (Σ, V, Q, S),
P ∗ := {α∗ → β ∗ : α → β ∈ P } and
P + := P ∗ ∪ {S → S ∗ }.
Show that for G+ := (Σ, V ∗ , P + , S), we have that L(G) = L(G+ ). [You may use without
proof that isomorphic grammars produce the same language.]
(c) Again, let G = (Σ, V, P, S) be an arbitrary grammar and use the notation
from (b). Let
Q+ := P + ∪ {S → ε} and
H + := (Σ, V ∗ , Q+ , S)
Paper 4, Section I
4J Automata and Formal Languages
(a) Let D = (Σ, Q, δ, q0 , F ) and D′ = (Σ, Q′ , δ ′ , q0′ , F ′ ) be two deterministic
automata. Define what it means that f is a homomorphism from D to D′ .
(b) Prove that if f is a homomorphism from D to D′ , then L(D) = L(D′ ).
(c) Define the following partial order on deterministic automata: we write D ⩽ D′ if
and only if there is an injective homomorphism from D to D′ . Check whether the following
statements are true or false. Justify your answers.
(e) Using part (d) or otherwise, show that L = {an bn : n > 0} is not 0-computable.
Paper 3, Section II
12J Automata and Formal Languages
(a) Define what it means for a language L ⊆ W to satisfy the regular pumping
lemma.
(b) Prove that every regular language satisfies the regular pumping lemma. [You
may assume that “regular” is equivalent to “accepted by a deterministic automaton”
without proof.]
(c) Let L be a finite language such that there is a w ∈ L with |w| = 100. Show that
there can be no deterministic automaton D with at most 100 states such that L = L(D).
(d) Let Σ = {0, 1, 2} and let L ⊆ {0, 1}∗ be an arbitrary language. Show that
b := {u2v : u ∈ W, v ∈ L} ∪ {0, 1}∗
L
Paper 2, Section I
8E Classical Dynamics
(a) In Lagrangian mechanics, explain what is meant by the generalised momentum
associated to a generalised coordinate q.
(b) What does it mean for a generalised coordinate to be ignorable? Show that the
generalised momentum associated to an ignorable coordinate is conserved. [You may state
the Euler–Lagrange equations without proof.]
(c) A certain system has generalised coordinates (q1 , q2 , q3 ) and Lagrangian
1 2 1 2
L= q̇1 + q̇22 + q̇32 − q1 + q22 + q32 − α (q1 q2 + q2 q3 + q3 q1 ) ,
2 2
where α is constant. Show that L is invariant under rotations around the (1, 1, 1) axis in
q-space. Hence find two conserved quantities.
Paper 3, Section I
8E Classical Dynamics
(a) Using spherical polar coordinates, (r, θ, ϕ), write down the Hamiltonian for a
particle of mass m moving in a spherically symmetric potential.
(b) Show that pϕ and p2θ +(p2ϕ / sin2 θ) are each conserved. Interpret them physically.
[You may state Hamilton’s equations without proof.]
(c) The particle executes circular motion in a plane through the origin inclined
at angle ψ to the plane θ = π/2. Evaluate pθ (θ) and show that it vanishes when
sin θ = ± cos ψ. Interpret this result physically.
|p|2
but with a free Hamiltonian, H = 2m , where qi and pi are generalised coordinates and
momenta, respectively.
Paper 2, Section II
14E Classical Dynamics
A mass m1 is suspended from a fixed point with coordinates (x, y, z) = (0, 0, 0) by
a spring with spring constant k1 . A second mass m2 is suspended from the first mass
by a spring with spring constant k2 . Each spring has natural length ℓ. The motion of
the masses is restricted to the (x, y)-plane, with gravity acting in the −y direction. The
position of the first mass is (x1 , y1 , 0) and the position of the second mass is (x2 , y2 , 0).
(a) Write down the Lagrangian of the system and hence determine the equations of
motion.
(b) Find the equilibrium position of each mass that has y2 < y1 < 0.
(c) For the remainder of this question suppose that the x-coordinate of mass m2 is
held fixed at its equilibrium value, and consider the case m1 = m2 = m and k1 = k2 = k.
One of the system’s normal modes has the first mass moving in the x direction with no
motion in the y direction and the other mass stationary. Show that this mode’s frequency
ω1 satisfies !
k 1 1
ω12 = 2 − 2mg − mg .
m kℓ + 1 kℓ + 1
Find the other normal modes and corresponding frequencies, showing that they are
independent of the strength of gravity.
(a) Write down expressions for the kinetic energy T and total angular momentum
L, and show that they are each conserved.
√(b) Suppose that√ I1 < I2 and I3 = I1 + I2 , and that initially ω2 (0) = 0 while
ω1 (0) I2 − I1 = ω3 (0) I2 + I1 . Show that subsequently
2 2T 2 I2 − I1
ω̇1 = − ω1 ω12 .
I2 I2 + I1
Find a binary Huffman code for X. What is its expected word length? [You do not
need to simplify the expression.]
Paper 2, Section I
3K Coding and Cryptography
(a) Show that Hamming’s original code is perfect.
(b) Consider the code obtained by using Hamming’s original code for the first 7 bits
and the final bit as a check digit, so that
x1 + x2 + · · · + x8 ≡ 0 (mod 2).
Find the minimum distance for this code. How many errors can it detect? How many
errors can it correct?
(c) Given a code of length n which corrects e errors, can you always construct a
code of length n + 1 which detects 2e + 1 errors? Give a brief justification of your answer.
Paper 3, Section I
3K Coding and Cryptography
In the following, equivalent definitions of the Reed–Muller code can be used without
justification.
(a) Let n = 2d (d ⩾ 1). For 0 ⩽ r ⩽ d, state and prove a formula for the rank of the
Reed–Muller code RM(d, r) of length n. State its minimum distance.
(b) Consider the Mariner 9 code RM(5,1). What is its information rate? What
proportion of errors can it correct in a single codeword? How do these two properties
compare to the Hamming code of length 31?
(c) Show that all but two codewords in RM(d, 1) have the same weight.
(a) (i) Consider a cryptosystem (K, M, C). Let e, d be the respective encryption
and decryption functions. Model the key and messages as independent
random variables K, M taking values in K, M, respectively and such that
M = d(C, K) ∈ M and C = e(K, M ) ∈ C. Show that H(M |C) ⩽ H(K|C).
(ii) Let M = C = A, where A is a finite alphabet. Suppose we send n messages
(M1 , . . . , Mn ) encrypted as (C1 , . . . , Cn ) using the same key. Define the
unicity distance. By making some reasonable assumptions, give a closed
formula for the unicity distance as a function of |K|, |A| and a certain constant.
(b) Suppose we model English text by a sequence of random variables (Xn )n⩾1 taking
values in A = {A, B, . . . , Z, space}. We define the entropy of English to be
Assuming HE exists, show that 0 ⩽ HE ⩽ log2 27. [You may assume Gibbs’
inequality.]
Paper 1, Section II
11K Coding and Cryptography
(a) Consider the use of a binary [n, m]-code to send one of m messages through a
binary symmetric channel (BSC) with error probability p, making n uses of the channel.
Define the following decoding rules: (1) ideal observer, (2) maximum likelihood, and (3)
minimum distance. Show that if all the messages are equally likely then (1) and (2) agree.
If p < 21 show that (2) and (3) agree.
1
(b) Show that a BSC with error probability p < 4 has non-zero operational capacity.
(c) State Shannon’s second coding theorem. Consider a discrete memoryless channel
with input X taking values over the alphabet {0, 1}. For a, b ∈ Z, let Z be a random
variable that is independent of X, taking values over the alphabet {a, b} with distribution
P(Z = a) = P(Z = b) = 12 . The output of the channel is Y = X + Z. What is the capacity
of this discrete memoryless channel? [Hint: The capacity depends on the value of b − a.]
(a) (i) Describe briefly the Rabin cryptosystem, including how to encrypt and decrypt
messages. Show that breaking the Rabin cryptosystem is essentially as
difficult as factoring the public modulus, N .
(ii) Criticise the following authentication procedure:
Alice chooses N as the public modulus for the Rabin cryptosystem. To be sure
you are in communication with Alice, you send her a “random item” r = m2
(mod N ). On receiving r, Alice proceeds to decode using her knowledge of
the factorisation of N , and finds a square root m1 of r. She returns m1 to
you and you check that r = m21 (mod N ).
(b) (i) Describe briefly the RSA cryptosystem with public modulus N .
A budget internet company decides to provide each of its customers with
their own RSA ciphers using a common modulus N . Customer j is given
the public key (N, ej ) and sent secretly their decrypting exponent dj . The
company then sends out the same message, suitably encrypted, to each of its
customers. You intercept two of these messages to customers i and j where
ei and ej are coprime. Explain how you would decipher the message.
You are one of the customers, and so also know your own decrypting exponent.
Can you decipher any message sent to another customer?
(ii) Explain why it might be a bad idea to use RSA with public modulus N = pq
with |p − q| small.
A user of RSA accidentally chooses the public modulus N to be a large prime
number. Explain why this system is not secure.
8πG
H2 = V (ϕ) ,
3c2
dV
3H ϕ̇ = −c2 ,
dϕ
where ti is the initial time with a(ti )=1 and ϕi =ϕ(ti ), with ϕi assumed to be large enough
to ensure inflationary expansion.
(b) By determining the Hubble parameter H, show that during inflation we have
1 2 c4 n2 V (ϕ)
ϕ̇ ≈ .
2c2 48πG ϕ2
Deduce the approximate value of ϕ = ϕend when inflation ends, that is, when the slow-roll
approximation breaks down. If n = 6, roughly estimate the initial value ϕi relative to ϕend
that would be required to solve the flatness problem of the standard cosmology.
where the deuterium binding energy is BD = (mp + mn − mD )c2 = 2.2 MeV, and
mp , mn and mD are the proton, neutron and deuterium masses, respectively. [Hint:
The degeneracy factor for Deuterium is gD = 4.]
(b) Now use fractional densities relative to the baryon number density nB (e.g.
Xp = np /nB ) to find an expression for XD /(Xp Xn ). In this case, replace nB = η nγ where
η is the baryon-to-photon ratio and the photon number is
16πζ(3)
nγ = (kB T )3 ,
(hc)3
where ζ is the Riemann zeta function. Briefly explain how the fractional density ratio
XD /(Xp Xn ) offers insight into the “deuterium bottleneck”, that is, the delay in forming
deuterium nuclei to temperatures well below the binding energy, kB T ≪ BD ?
8πG 2 4πG
H2 + kc2 = ρa , H′ = − (ρ + 3P )a2 ,
3c2 3c2
where ρ is the energy density, P is the pressure, k is the curvature and primes denote
differentiation with respect to conformal time τ (defined by dτ = dt/a(t)). Here, a is the
scale factor and H is the conformal Hubble parameter defined by H ≡ a′ /a.
(a) Assume that the universe is filled with a single-component fluid with equation
of state P = wρ, where w is a constant. By introducing the density parameter
Ω = 8πGa2 ρ/(3c2 H2 ), show that the two evolution equations can be recast as
kc2
= Ω − 1, 2H′ + (1 + 3w) H2 + kc2 = 0 .
H2
Hence, find the following evolution equation for the density parameter,
Ω′ = (1 + 3w)HΩ(Ω − 1) . (†)
(b) Using the time evolution of Ω from equation (†), qualitatively describe the
flatness problem of an expanding universe (H > 0) for models with an equation of state
parameter in the range 0 ⩽ w ⩽ 1. In particular, roughly sketch the time evolution of Ω
in a radiation-filled universe with w = 1/3 taking initial values for both Ω < 1 and Ω > 1.
(c) Briefly discuss how the flatness problem can be alleviated by an early epoch of
inflation during which w ≈ −1.
dE = T dS − P dV .
(a) By substituting the Stefan–Boltzmann law, show that the entropy differential
becomes
16σ 3
dS = T dV + 3T 2 V dT ,
3c
which should be integrated to find an expression for the photon entropy density s = S/V .
(b) If the interaction rate Γ maintaining the photons in equilibrium is much greater
than the Hubble expansion rate H (i.e. Γ ≫ H), briefly give the key reason why the
photon number and entropy are conserved as the Universe expands (provided the effective
number of degrees of freedom g∗ of particle species in equilibrium also does not change).
Why does the photon temperature fall as T ∝ 1/a, where a is the scale factor?
(c) Electrons and positrons annihilate and fall out of equilibrium after neutrino
decoupling at around kB T ≈ 1 MeV. By counting the effective number of degrees of freedom
g∗ in equilibrium before and after this process, provide a brief explanation for why the
photon temperature Tγ and neutrino temperature Tν are related today by
1/3
Tν 4
= .
Tγ 11
where k is a constant.
(b) Now differentiate the Friedmann equation (†) and substitute the continuity
equation (⋆) to find the acceleration equation for ä/a in terms of ρ, P and Λ. Briefly note
two of the shortcomings of this Newtonian analysis.
(c) Consider a flat (k=0) universe with a positive cosmological constant Λ > 0 that
is filled with radiation pressure PR = ρR /3, measured to have energy density ρR (t0 ) = ρR0
at given time t = t0 . Use the Friedmann equation (†) to show that the Hubble parameter
H = ȧ/a can be expressed as
1 8πG ρR0
H 2 = H02 ΩR0 a−4 + Λc2 , where ΩR0 ≡ ,
3 3c2 H02
with H(t0 ) =H0 and a(t0 ) = 1. By considering the substitution b = a2 (or otherwise) find
the solution for the scale factor
where α and β are constants you should determine in terms of H0 and ΩR0 . [You may
assume that the universe started with a big bang.]
Show that the scale factor a(t) gives anticipated results at early and late times.
Estimate the transition time tΛ that separates the decelerating and accelerating epochs.
R √
[Hint: dx/ 1 + κ2 x2 ) = (1/κ) sinh−1 (κx) + const, where κ > 0 is a constant.]
δ̇ + ∇ · [v(1 + δ)] = 0 ,
ȧ c2
v̇ + 2 v + (v · ∇)v = −∇ϕ − ∇P ,
a ρ
4πG
∇2 ϕ = 2 ρ̄ δ ,
c
where P (x, t) is the pressure and ϕ the perturbation of the Newtonian potential. Here, we
will neglect rotational modes and assume the velocity v is solely compressional, that is,
the divergence is the only non-vanishing part and θ ≡ ∇ · v. We also assume that pressure
effects can be described using the sound speed c2s defined by c2s ≡ c2 (dP/dρ).
(a) Show that for an equation of state given by P = P (ρ), the pressure term
in the linearised Euler equation can be written as (c2 /ρ)∇P ≈ c2s ∇δ. Linearise both
the continuity equation and the divergence of the Euler equation (using θ = ∇ · v) and
substitute the Poisson equation in the latter. Transforming to Fourier space with comoving
wavemodes k (i.e. ∇ → ik/a), combine these to obtain the evolution equation for density
perturbations in an expanding universe,
ȧ 4πG c2s k 2
δ̈ + 2 δ̇ − ρ̄ − 2 δ = 0, (†)
a c2 a
where k = |k|. Briefly explain the qualitative implications of each term in this equation.
Define the Jeans length λJ and discuss its significance.
Given k̃J = 5/(6L0 ), describe these solutions at late times in the asymptotic limits where
k ≪ k̃J and where k ≫ k̃J . What would be the consequence of choosing α such that
L0 ≈ 50 Mpc today?
ϕ : Ω → R3
(x, y) 7→ (x, y, h(x, y))
ϕt : Ω → R3
(x, y) 7→ ϕ(x, y) + tN (x, y)
and assume that there is ϵ > 0 so that St := ϕt (Ω) is a smooth surface for any t ∈ (−ϵ, ϵ).
Prove that the second fundamental form of S satisfies
1d
S
IIϕ(x,y) =− IϕStt(x,y)
2 dt t=0
S
at any point (x, y) ∈ Ω. Calculate IIϕ(x,y) in terms of h.
(d) Derive a differential equation in h that characterises when the surface S is
minimal.
(e) Calculate the area of S in terms of the height function h. Assume that for any
η : Ω → R smooth and compactly supported in Ω, the area Aη (t) of
(d) Consider the curve γ ϵ := (α + ϵβ)/|α + ϵβ| for small ϵ. You may assume that,
RL
if the value ϵ = 0 is a critical point of the enclosed area, then 0 κ−1
g φ̇ = 0. Deduce from
this equation that area-maximising curves have constant geodesic curvature.
(e) Prove that a curve α on S2 with constant geodesic curvature is planar by showing
that the vector e(s) := α(s) × α̇(s) + κg α(s) is constant.
(f) Deduce that, if a curve on S2 of length L encloses an area A, then L2 ⩾ A(4π−A)
(with the convention that we always choose the smaller of the two areas enclosed by the
curve).
Paper 3, Section II
25J Differential Geometry
Let n ⩾ 1 and 1 ⩽ k ⩽ N be integers. Let Id denote the identity matrix.
(a) State the definition of a k-dimensional smooth manifold X ⊂ RN . Define the
tangent space Tx X at a point x ∈ X in terms of a parametrisation about x, and prove
that the tangent space is independent of the parametrisation.
2
(b) Let GL(n) ⊂ Rn be the set of real n × n invertible matrices. Prove that it is a
manifold, give its dimension, and compute TId GL(n).
2
(c) Let SL(n) ⊂ Rn be the set of real n × n matrices with determinant 1. Prove
that it is a manifold, give its dimension, and compute TId SL(n).
2
(d) Let SO(n) ⊂ Rn be the set of real n × n matrices with determinant 1 and
orthonormal column vectors. Prove that it is a manifold, give its dimension, and compute
TId SO(n).
x + ϵφ̃(x)
ψϵ (x) := √ .
1 + ϵ2
Prove that, for ϵ small enough, ψϵ is a diffeomorphism from A to A with det Dψϵ > 0.
R
(c) Prove A det(Id + ϵDφ̃) = vol(A)(1 + ϵ2 )(n+1)/2 , where vol(A) is the volume of
A.
(d) Deduce that, when n is even, there is no smooth map φ : Sn → Sn such that
φ(x)⊥x for all x ∈ Sn .
(i) Show that F also has periodic orbits of period n for all positive integers n.
(ii) Explain briefly why F must have at least four distinct 7-cycles.
(iii) How many distinct 8-cycles must F have?
[Relevant theorems that you use from the course should be stated clearly.]
ẋ = y − x + ax3 ,
ẏ = rx − y − zy,
ż = −z + xy,
(i) Show that the fixed point at the origin of the system is non-hyperbolic at
r = 1.
(ii) Find the stable, unstable and centre subspaces of the linearised system of
the fixed point at the origin at r = 1.
(iv) Show that the evolution equation on the centre manifold is of the form
a − 1 3 (3a − 1)(a + 3) 5
v̇ = v + v + ...
2 8
Z T
exp ∇ · f (X(t)) dt .
0
(b) Using the energy-balance method for nearly Hamiltonian systems, find a
condition on the parameter a for a limit cycle to exist in the system given by
ẋ = y,
ẏ = −x + ϵ(1 − x2 + ay 2 )y,
(ii) Explain why every periodic orbit must enclose at least one fixed point.
ẋ = −x + ay + x2 y,
ẏ = b − ay − x2 y,
(i) Find the fixed point of the system. On a sketch of the (a, b)-plane (for the
quadrant where a, b > 0), show where the fixed point is stable.
(ii) A closed region D of the (x, y)-plane is given by the filled polygon with
vertices at (0, 0), (0, b/a), (b, b/a), (x∗ , y ∗ ) and (x∗ , 0), where (x∗ , y ∗ ) is
such that ẋ = 0 at that point and the line segment between (b, b/a) and
(x∗ , y ∗ ) has gradient −1.
Show that trajectories do not leave D.
[There is no need to determine (x∗ , y ∗ ) explicitly, but you may use that
b < x∗ .]
(iii) Use your results above to give conditions on the parameters a and b for the
system to have a periodic orbit.
∂µ F µν − m2 Aν = −µ0 J ν ,
□A − ∇α − m2 A = −µ0 J ,
2
where α = 1 ∂ϕ
c2 ∂t
+ ∇ · A and □ = − c12 ∂t
∂ 2
2 + ∇ is the wave operator.
(e) For a point charge q at rest, the 4-current density is J 0 = cqδ(x) and J = 0,
where δ denotes the 3-dimensional δ function. By applying a Fourier transform in the
spatial coordinates to the equation of motion (†), show that for m > 0, a time-independent
field solution for a point charge at r = 0 is given by
exp(−mr)
ϕ=λ , A = 0,
r
where r = |x| and λ is a constant you do not need to determine. [Hint: You may use
without proof that the inverse Fourier transform of 1/(|k|2 + m2 ) is
Z
eik·x d3 k e−m|x|
= C ,
|k|2 + m2 (2π)3 |x|
ρ ∂B
∇·E = , ∇×E =− ,
ϵ0 ∂t
∂E
∇ · B = 0, ∇ × B = µ0 J + µ0 ϵ0 .
∂t
The energy and momentum density of the electromagnetic field are defined by
1 2 B2
E= ϵ0 E + and g = ε0 E × B .
2 µ0
(a) Use the Maxwell equations to show that g obeys the local conservation law
with
σij = ã Ei Ej + b̃ E 2 δij + c̃ Bi Bj + d˜B 2 δij ,
where δij is the Kronecker delta and ã, b̃, c̃, d˜ are constants you should determine. [Hint:
A vector field a satisfies (∇ × a) × a = (a · ∇)a − 12 ∇(a2 ) ]
(b) Provide a brief physical interpretation of the tensor σij and, in particular, its
diagonal and off-diagonal components.
(c) The electric field of a homogeneously charged sphere with total charge q and
radius R centered on the origin is given by
(
q r
3 for r ⩽ R
E = E(r)er with E(r) = 4πϵ 0 R
q 1
4πϵ0 r2 for r > R ,
where r = |x| and er is the unit vector in the radial direction. The magnetic field B is
zero. Calculate the energy density E and, thus, the total energy Wem contained in the
electromagnetic field.
Under a Lorentzp transformation with velocity v = (0, 0, v) in the z′ direction and
2 2 ′
Lorentz factor γ = 1/ 1 − (v /c ), the electromagnetic field changes to E and B given
by
E ′|| = E || , E ′⊥ = γ(E ⊥ + v × B) ,
v E
B ′|| = B || , B ′⊥ = γ B ⊥ − × ,
c c
where the subscripts || and ⊥, respectively, denote field components parallel and perpendic-
ular to v. Compute the linear momentum P = (0, 0, Pz ) contained in the electromagnetic
field of a homogeneously charged sphere moving with v in the slow-velocity limit, i.e. ig-
noring terms of order (v/c)2 or higher.
R π [Hint: Vector fields a, b and c satisfy the relation
a × (b × c) = (a · c)b − (a · b)c and 0 sin θdθ = 43 .]
3
(b) Starting from the microscopic Maxwell equations for the electric and magnetic
fields E, B,
ρ ∂B
∇·E = , ∇×E =− ,
ϵ0 ∂t
∂E
∇ · B = 0, ∇ × B = µ0 J + µ0 ϵ0 ,
∂t
derive the macroscopic Maxwell equations for a dielectric medium in terms of the electric
displacement D and the magnetising field H, which you should define, as well as the free
charge and current densities.
(c) Consider a spherical shell with radial extent R1 < r < R2 , that consists of a
linear magnetic medium with constant permeability µ such that B = µH. This shell is
placed inside a magnetic field that approaches B = B0 ez , B0 = const, as r → ∞. The
electric field, the polarization, and the free charge and current densities are zero. Using
the ansatz B = ∇ψ, show that a solution for the magnetic field in all space is given by
(a1 r + b1 /r2 ) cos θ for r < R1
ψ(r, θ) = (a2 r + b2 /r2 ) cos θ for R1 < r < R2
(a r + b /r2 ) cos θ
3 3 for r > R2 .
Using the boundary conditions at infinity, the origin, r = R1 and r = R2 , derive six
conditions for the free parameters a1 , b1 , a2 , b2 , a3 and b3 . Express the parameters b1 , a2 ,
b2 and a3 in terms of B0 , µ, R1 , R2 and a1 . Briefly describe how you would calculate the
remaining parameters a1 and b3 (you do not need to compute a1 and b3 ).
[Hint: In polar coordinates, the Laplace operator and the derivative in the Cartesian
z direction are
2 1 ∂2 1 ∂ ∂ 1 ∂2
∇ f = (rf ) + sin θ f + f,
r ∂r2 r2 sin θ ∂θ ∂θ r2 sin2 θ ∂ϕ2
∂ ∂ sin θ ∂
f = cos θ f− f .]
∂z ∂r r ∂θ
Paper 2, Section II
39C Fluid Dynamics II
A two-dimensional incompressible Stokes flow has stream function ψ such that the
velocity u = ∇×(ψk), where k is the unit vector normal to the plane of the flow. Show
that
∇4 ψ = 0.
1 ∂ψ ∂ψ
u= er − eθ .
r ∂θ ∂r
Given that the stream function has the form ψ = r2 f (θ), determine the rate-of-strain
tensor in terms of f and its derivatives. Hence write down the corresponding deviatoric
stress tensor for a fluid of dynamic viscosity µ.
Fluid with dynamic viscosity µ fills the two-dimensional region −α < θ < 0, r > 0,
where α > 0 is a constant. The boundary θ = −α is rigid, while a tangential stress S
is applied to the horizontal surface θ = 0. Given that the stream function has the form
ψ = r2 f (θ), write down the boundary conditions that apply to f (θ). Hence, determine
f (θ) and show that the surface velocity
Sr 1 − cos 2α − α sin 2α
U (r) = u(r, 0) = .
µ sin 2α − 2α cos 2α
2 1 ∂ ∂ψ 1 ∂2ψ
[Hint: In plane polar coordinates, ∇ ψ = r + 2 2 .]
r ∂r ∂r r ∂θ
Paper 4, Section II
38C Fluid Dynamics II
Given a material interface r = R(θ, t) between two fluid regions in plane polar
coordinates (r, θ), explain why
∂R v ∂R
+ = u,
∂t r ∂θ
where uer + veθ is the fluid velocity at the interface.
An ocean gyre is modelled as a two-dimensional, circular patch of water of radius a
in solid-body motion with angular velocity ω, surrounded by stationary water. Consider
small, sinusoidal perturbations to the edge of the gyre r = a + η(θ, t), where η =
ϵ exp(ikθ + σt) and ϵ ≪ ka. Assuming potential flow, determine the relationship between
the growth rate σ and the wave number k.
Briefly describe the subsequent motion of the gyre. Do disturbances propagate
upstream or downstream relative to the original motion of the gyre?
2 1 ∂ ∂ϕ 1 ∂2ϕ
[Hint: In plane-polar coordinates, ∇ ϕ = r + 2 2 .]
r ∂r ∂r r ∂θ
Paper 2, Section I
7D Further Complex Methods
Consider the differential equation
d3 y(x)
x + 2y(x) = 0 (†)
dx3
on the domain x ∈ (0, ∞).
(a) Write Z
y(z) = ezt f (t) dt,
γ
where γ is a contour in the complex plane, and substitute this expression for y into the
differential equation (†). Explain how the resulting integral equation can be solved by
finding an appropriate function f (t) and contour γ. Determine this function f (t) and
clearly state any required conditions on γ.
(b) Express the solution y(x) in integral form. [You do not have to evaluate this integral,
but you should simplify it as far as possible.] [Hint: Consider a subset of the real axis for
your choice of the contour γ.]
1. F (z) is analytic on H.
4. F (1) = 1.
Ly(t) = f (t),
where f (t) is the input, y(t) is the output and L is a linear operator. Define what it means
for the system to be causal and stable.
(b) Consider a linear ordinary differential equation for t > 0,
(i) Using a Laplace transform, show that the transfer function of the linear
1
system (†) is G(s) = αs+1 . Determine for which values of α the system is
stable.
Can there exist an elliptic function with a single pole of order one in a fundamental cell?
(d) If h is a meromorphic function on and inside a simple closed clockwise contour
γ and h has no zeros or poles on γ, then
I ′
1 h (z)
dz = P − Z ,
2πi γ h(z)
where P and Z denote respectively the number of poles and zeros, counting multiplicities,
of h(z) inside the contour γ. Using this relation, show that a non-constant elliptic function
f takes each value the same number of times in a cell, counting multiplicities.
(e) An example of an elliptic function is the Weierstrass function
1 X 1 1
P(z) = 2 + − ,
z (z − wm,n )2 wm,n2
(m,n)
X ∞
1
P(z) = 2 + a2k z 2k .
z
k=0
Paper 2, Section II
18H Galois Theory
(a) Let L/K be a field extension. Explain what it means to say that
Which pairs of these properties together imply the third? In each case give a proof or
counterexample.
(b) Let L be the splitting field of f (X) = X 3 − X − 1 over Q. Compute√Gal(L/Q).
Show that L has a unique quadratic subfield, and write it in the form √ Q( d) for d a
squarefree integer. Show also that if α is a root of f then L = Q(α + d).
Paper 3, Section II
18H Galois Theory
(a) Let M/L1 /K and M/L2 /K be finite extensions of fields. Define the composite
L1 L2 . Show that if L1 /K is Galois then L1 L2 /L2 is Galois, and that there is an injective
group homomorphism Gal(L1 L2 /L2 ) → Gal(L1 /K).
(b) Let K be a field of characteristic not dividing n, and with algebraic closure K.
Let ζn ∈ K be a primitive nth root of unity. Prove that [K(ζn ) : K] divides ϕ(n), and
that equality holds if K = Q.
(c) Let L1 = Q(ζ9 ) and L2 = Q(ζ15 ). Write each of the fields L1 ∩ L2 and L1 L2 in
the form Q(ζn ) for suitable n. Justify your answers.
Consider a massive test particle moving in the Schwarzschild metric of a black hole
with mass m (in units with c = G = 1):
2m 2 2m −1 2
ds2 = − 1 − dt + 1 − dr + r2 ( dθ2 + sin2 θ dϕ2 ) .
r r
(a) Assuming that the motion lies in the equatorial plane θ = π/2, justify briefly
why
1 2 m h2 mh2
h = r2 ϕ̇ and ṙ − + 2− 3
2 r 2r r
are constants of the motion, where dots denote derivatives with respect to the proper time
of the particle.
(b) For a circular orbit with a fixed value of r, determine h and hence deduce that
(i) r > 3m and (ii) (dϕ/dt)2 = m/r3 .
(c) Now consider a nearly circular orbit with shape given by u(ϕ) = 1/r. Let prime
denote differentiation with respect to ϕ, so that u′ = du/dϕ. Given that ṙ = −hu′ , and
assuming if needed that u′ ̸= 0, show that
m
u′′ + u = + 3mu2 .
h2
For m/h ≪ 1, this equation has an approximate solution of the form
m
u = (1 + α) + A cos[(1 + β)ϕ] ,
h2
where the constant A obeys |A| ≪ 1 but is otherwise arbitrary. The constants α and β
are small for m/h ≪ 1. Verify this solution by working to first order in A and determining
the constants α and β to leading non-trivial order in m/h.
Comment briefly on the significance of your result for β.
(a) Define the Einstein tensor Gµν in terms of the Riemann tensor Rα βµν and use
the Bianchi identity ∇ρ R α βµν + ∇µ R α βνρ + ∇ν R α βρµ = 0 to show that ∇ν Gµ ν = 0.
Comment briefly on the significance of this result for consistency of the Einstein equations
(including a cosmological constant).
(b) For a universe described by the line element
the Einstein tensor is diagonal with Gt t = −3ȧ2 /a2 and Gx x = −2ä/a − ȧ2 /a2 , where dots
denote differentiation with respect to t. Verify by direct computation that ∇ν Gµ ν = 0,
justifying the steps that you make and computing any metric connection components Γνµρ
that you may need.
Solve the vacuum Einstein equations with a cosmological constant Λ > 0 to obtain
a result for a(t) that corresponds to an expanding universe.
(a) Consider a curve xα (λ) with tangent vector T α = dxα /dλ, in a spacetime with
metric gµν . For any vector fields U α and W α , show that, on the curve,
d
( gµν U µ W ν ) = (∇T U )α Wα + Uα (∇T W )α ,
dλ
where ∇V = V α ∇α denotes the covariant directional derivative along a vector V α .
(b) Now consider a one-parameter family of geodesics defined by the functions with
two arguments, xα (τ, σ). Fixing a value of σ in these functions gives a timelike geodesic
parametrised by proper time τ , with tangent vector T α = ∂xα /∂τ satisfying T α Tα = −1,
while fixing a value of τ gives a curve parametrised by σ, with tangent vector S α = ∂xα /∂σ.
Show that the commutator [T, S]α = (∇T S)α − (∇S T )α = 0 and hence derive the
equation of geodesic deviation in the form
(∇T ∇T S)α + E α β S β = 0 ,
(∇T ∇T S)α − KS α = 0
(a) For a spacetime with metric gαβ , write down an explicit expression for the
metric-preserving Levi-Civita connection Γβαγ .
For a spacetime that is nearly flat, the metric can be expressed in the form
where ηαβ is a flat metric with constant components (though not necessarily diagonal in
the coordinates used) and the components hαβ and their derivatives are small. Show that
to leading order in these small quantities,
2Rαβ = A hα γ , γβ + B hβ γ , γα + C hγ γ , αβ + D hαβ, γρ η γρ ,
for constants A, B, C, D which you should determine. Indices are raised and lowered using
ηαβ .
(b) Consider the following metric
where H(u, x, y) is a smooth function that is not necessarily small. You may assume that
the only nonzero connection coefficients are Γxvu , Γyvu , Γuxu , Γuyu , Γuvu , and the coefficients
related to these by symmetry. You may also assume that Ruu is the only coefficient of the
Ricci tensor that is not trivially zero. Compute Ruu and hence obtain the Ricci scalar.
Show that the full nonlinear vacuum field equations, with no cosmological constant, reduce
to a partial differential equation for H, that you should determine.
[You may assume in both parts of the question that
Paper 2, Section II
17I Graph Theory
State and prove Turán’s theorem.
A rhombus is the graph formed by two triangles sharing an edge. Prove that if G is
a graph on n ⩾ 4 vertices that has more edges than T2 (n), then G contains a rhombus.
Find a graph G on 6 vertices such that e(G) = e(T2 (6)) and G does not contain a
rhombus, but G is not isomorphic to T2 (6).
Paper 3, Section II
17I Graph Theory
(a) Show that every graph of average degree at least d contains a subgraph of
minimum degree at least d/2.
(b) Let δ, g ⩾ 3 be fixed. Using random graph methods, or otherwise, show that
there exists a graph G on at most n = (100δ)g vertices such that the minimum degree of
G is at least δ and the girth of G is at least g.
are integers.
What are the eigenvalues (with their multiplicities) of the Petersen graph (shown
below)?
Show that the set of edges of K10 cannot be partitioned into the edges of three
copies of the Petersen graph. [Hint: Suppose that it can be, and that the three Petersen
graphs have adjacency matrices A, B and C. You may wish to consider an appropriate
common eigenvector of A and B.]
Show that the equation ut = iuxx has a formulation as the consistency condition
for the following pair of equations for ψ = ψ(x, t) ∈ C,
ψt + ik 2 ψ = iux − ku,
ψx − ikψ = u.
2
By means of the integrating factor e−ikx+ik t , or otherwise, deduce that
Z t
ik2 t 2
û(k, t)e − û0 (k) = eik τ [ku(0, τ ) − iux (0, τ )] dτ ,
0
where
Z ∞ Z ∞
−ikx
û(k, t) = e u(x, t)dx, and û0 (k) = e−ikx u0 (x)dx .
0 0
Hence, by considering also û(−k, t), find a function G = G(k, τ ) such that
Z +∞
1 2
u(x, t) = e−ik t+ikx [û0 (k) − û0 (−k)] dk
2π −∞
Z Z
1 +∞ t −ik2 (t−τ )+ikx
+ e G(k, τ )dτ dk .
π −∞ 0
[You may assume that u is smooth and rapidly decreasing so that û(k, t) is holomorphic
for Im{k} < 0, and satisfies lim|k|→∞ û(k, t) = 0.]
(b) Consider now the case of periodic boundary conditions for the KdV equation, so
that at each time t the unknown u is a real-valued function satisfying u(x+2π, t) = u(x, t).
At each fixed time t, introduce a basis φ+ , φ− of solutions to the scattering equation
Lϕ = k 2 ϕ (for k 2 > 0) determined by the initial conditions at x = 0,
where a(t) and b(t) are functions of time and a, b denote the complex
conjugates, such that
φ+ (x + 2π) φ+ (x)
= T̂ .
φ− (x + 2π) φ− (x)
(ii) Now as t varies let u evolve in time according to the KdV equation. Show
that there exists a matrix
λ µ
Λ= ,
µ λ
pr (N ) V = V1 ∂x + ϕ ∂u + ϕ1 ∂u′ + · · · + ϕN ∂u(N )
of the vector field V = V1 (x, u)∂x + ϕ(x, u)∂u . Give the inductive formula for the
computation of the ϕj = ϕj (x, u, u′ , . . . , u(j) ) which determines the prolongation.
For the special case in which V1 = f (x) depends only on x, compute the third
prolongation, and show that
ϕ3 = ϕxxx +u′ (αϕxxu −f ′′′ )+3(u′ )2 ϕxuu +(u′ )3 ϕuuu +β(ϕxu −f ′′ +u′ ϕuu )u′′ +(ϕu −3f ′ )u′′′ ,
d3 u 1
3
= 3
dx u
that are generated by a vector field of the form V = f (x)∂x + ϕ(x, u)∂u .
Paper 2, Section II
22G Linear Analysis
State and prove the Baire Category Theorem.
Let X be a Banach space, and let S be a non-empty subset of X that is closed, convex
and symmetric (S symmetric means x ∈ S implies −x ∈ S). Show that if ∪∞ n=1 nS = X
then S is a neighbourhood of the origin.
Give an example to show that the condition that S is convex cannot be omitted.
Paper 3, Section II
21G Linear Analysis
State and prove the Stone–Weierstrass theorem. [You may assume that the function
x1/2 is uniformly approximable by polynomials on [0, 1].]
Let C(R) denote the space of all bounded continuous functions from R to R,
equipped with the uniform norm. Explain briefly why C(R) is a Banach space.
If A is a subalgebra of C(R) that contains the constants and separates the points,
must A be dense in C(R)? Justify your answer.
fv = {(x, y) ∈ X × Y : v(px,y ) = 1} .
For each of the following statements, either write down a set S ⊂ L that makes the
statement true or prove that no such set S exists.
(i) {fv : v is a model of S} is the set of all injective functions from subsets of
X to Y .
X (0) = X
′
X (α+1) = X (α)
\
X (λ) = X (α) (for non-zero limit ordinal λ)
α<λ
Show that if X ̸= ∅, then X ′ ̸= X. Deduce that X (α) = ∅ for some α. The least such α is
the index of X.
Show that if ξ is an ordinal, then
Paper 3, Section II
16I Logic and Set Theory
(a) State Zorn’s Lemma, the Axiom of Choice and the Well-ordering Principle, and
prove that they are equivalent.
(b) State Gödel’s Completeness Theorem and the Compactness Theorem for First-
order Logic.
Let T0 , T1 , . . . , Tn be first-order theories in some language L that partition the
collection of all L-structures in the sense that every L-structure is a model of exactly
one Ti . Show that each Ti is finitely axiomatisable.
ℵα + ℵβ = ℵα · ℵβ = ℵβ
for ordinals α ⩽ β.
Prove that the following sentence is a theorem of ZFC, where t ≡ x is taken to mean
‘t and x have the same cardinality.’
Paper 2, Section I
6A Mathematical Biology
The model of a viral infection in a population is given by the system
dX
= µN − βXY − µX,
dt
dY
= βXY − (µ + ν)Y,
dt
dZ
= νY − µZ,
dt
where µ, β and ν are positive constants and X, Y , and Z are respectively the number of
susceptible, infected and immune individuals in a population of size N , independent of t,
where N = X + Y + Z.
(a) Interpret the biological meaning of each of the parameters µ, β and ν.
(b) Show that there is a critical population size Nc (µ, β, ν) such that if N < Nc
there is no steady state with the infection maintained in the population. Show that in
this case the numbers of infected and immune individuals decrease to zero for all possible
initial conditions.
(c) Show that for N > Nc there is a steady state (X, Y, Z) = (X ∗ , Y ∗ , Z ∗ ) with
0 < X ∗ , Y ∗ , Z ∗ < N . Show that this steady state is stable.
Paper 4, Section I
6A Mathematical Biology
A discrete-time model of alien cell population dynamics considers the coupled
dynamics of immature and mature cells. At each time step, a proportion λ of the immature
cells mature and a proportion µ of mature cells divides. When a mature cell divides it is
replaced by four new immature cells. Finally, a proportion k of mature cells die at each
time step.
(a) Explain briefly how this model may be represented by the equations
where 0 ⩽ µ, λ, k ⩽ 1 and µ + k ⩽ 1.
(b) What is the expected number of offspring per cell?
(c) By considering the total population of cells, show that population growth is only
possible if 3µ > k. How does this relate to your answer to part (b)?
(d) At time T , the population is treated with a chemical that completely stops cells
from maturing for all t ⩾ T , but otherwise has no direct effects. Explain what will happen
to the population afterwards. In terms of aT and bT , what will be the total number of
cells in the long term?
∂C ∂2C
= D 2 + µC ,
∂t ∂x
where D and µ are positive constants.
Initially, the bacteria are at a uniform concentration C0 for 0 < x < L.
(a) Suppose that the ends of the channel are open so that bacteria may diffuse out
of the ends.
(ii) What is the flux of bacteria out of the channel? Comment on this flux at
t = 0.
(iii) Show that the population will grow if µ > µc , where µc (D, L) should be
given. Give a brief explanation for the dependence of µc on D and L.
(b) The experiment is run again with the x = L end of the channel closed (x = 0
remains open). Again initially C(x, t) = C0 . What is the condition now for population
growth in the long term? Comment briefly on how this compares to the condition for
growth in part (a).
(c) The experiment is run once more with the channel ends both open, but now the
per capita growth rate is a function of time (as the experimental conditions fluctuate each
day), given by µ(t) = µ0 + µ1 cos(t). Find the condition for population growth in the long
term. Give a brief interpretation of this result.
∂u ∂2u
= D1 2 + u − uv + u2 ,
∂t ∂x
∂v ∂2v
= D2 2 + αu2 − βuv,
∂t ∂x
where α, β, D1 , D2 > 0.
(a) Find conditions on α and β for the spatially homogeneous system to have a
stable stationary solution with u > 0 and v > 0. Sketch this region in the α–β plane (for
the quadrant with α, β > 0).
(b) Consider spatial perturbations of the form cos(kx) about the solution found in
part (a), and set λ = D1 /D2 . Find conditions for the system to be unstable. For fixed β,
sketch in the λ–α plane the region where spatial instability is possible for some k ∈ R.
(c) Consider a general system of the form
∂u ∂2u
= D1 2 + f (u, v),
∂t ∂x
∂v ∂2v
= D2 2 + g(u, v).
∂t ∂x
(0) (1)
(b) A new painkilling drug is tested on n patients. Let Yi and Yi be the
pain levels, on a scale from 0 (no pain) to M > 0 (maximum pain), of the ith patient,
before and after taking the painkiller respectively. Suppose the vector Xi ∈ Rp records
p additional
characteristics
of the ith patient, such as their age, weight, height, etc. We
(0) (1)
treat Yi , Yi , Xi ∈ [0, M ]2 × Rp for i = 1, . . . , n as independent copies of a random
triple Y (0) , Y (1) , X . Let A and H be defined as in part (a) (iii) above. We wish to
determine a region A ∈ A where if X ∈ A, we expect the drug to be effective. To this
end, let h∗ and ĥ minimise
h i 1 X (1)
n
(1) (0) (0)
Q(h) := E Y −Y h(X) and Q̂(h) := Yi − Yi h(Xi )
n
i=1
respectively, over h ∈ H.
(i) Using any results from the course that you need, show that
(ii) Using any results from the course that you need, show that
r
2p log(n + 1)
EQ(ĥ) ⩽ Q(h ) + 2M
∗
.
n
Now suppose further that Var(Y | X = x) is bounded from above by σ 2 for all
x∈ Rp . Show that
h 2i σ2J
E T̃ (X) − E T̃ (X) | X, X1:n ⩽ .
n
[Hint: You may use without proof, the fact that if N ∼ Binomial(n, q), for success
probability q ∈ (0, 1], then E[1/(N + 1)] ⩽ 1/(nq).]
(d) Now take ϕ to be the hinge loss. Writing ki ∈ Rn for the ith column of XX ⊤ ,
show that
n
1X
ki ti = 2λXX ⊤ α̂,
n
i=1
where ti is related to yi x⊤i β̂ in a way you should specify. Hence conclude that whenever
⊤
yi xi β̂ > 1, then α̂i = 0.
where d is the degree, ad is the leading coefficient and α1 , . . . , αd are the roots of f in C.
Suppose f ∈ Z[x] is irreducible, M (f ) = 2, and f has a real root α1 > 1.
Prove that α1 is an algebraic integer and |NQ(α1 )/Q (α1 )| = 2. [Hint: Consider the
number |NQ(α1 )/Q (α1 )|/2, and show that it is a rational integer.]
Paper 2, Section II
20F Number Fields
State Dirichlet’s unit theorem.
√
Let K = Q( 5), and determine the units in OK . [You may use without proof the
description of OK , as long as you state it clearly.]
√
For K = Q( 5), what are the possible degrees of extensions L/K of number fields
× ×
such that 1 < |OL /OK | < ∞? Give an example for each possible degree.
Paper 4, Section II
20F Number Fields
Let K be a number field.
Define the norm N (I) of an ideal I in OK .
Let d = [K : Q]. Define the discriminant of a tuple (α1 , . . . , αd ) ∈ K. Define the
discriminant of K.
State a formula for the norm of an ideal in terms of discriminants without proof.
Prove that N (αOK ) = |NK/Q (α)| for all α ∈ OK .
Now let L/K be an extension of number fields, and let P ⊂ OK and Q ⊂ OL be
non-zero prime ideals.
(a) Prove that Q|P OL if and only if P = Q ∩ OK .
(b) Suppose that P = Q ∩ OK and [OK : P ] = [OL : Q]. Prove that if Q = αOL for
some α ∈ OL , then P = NL/K (α)OK .
Paper 2, Section I
1F Number Theory
Let d be a positive integer.
Define what it means for a positive definite binary quadratic form f (x, y) =
ax2+ bxy + cy 2 to be reduced. If d is congruent to 0 or 3 mod 4, define the class number
h(−d).
Show that if d is odd and has k distinct prime factors, then h(−4d) ⩾ 2k−1 .
Give an example with d > 1 to show that the inequality h(−4d) ⩾ 2k−1 can be
strict.
Paper 3, Section I
1F Number Theory
Let N be an odd positive integer, and let a be an integer.
Define the Jacobi symbol Na , and write down a formula for 2
N
.
State
the law of quadratic reciprocity for the Jacobi symbol, and use it to com-
3
pute N in terms of the value of N modulo 12.
Show that if d is a positive integer suchthat d ≡ 0 or 3 mod 4, and a, b are positive
odd integers such that a ≡ b mod d, then −d −d
a = b .
Paper 3, Section II
11F Number Theory
Let k ∈ N.
Let p be an odd prime and Hk = {a + pk Z | a ≡ 1 mod p}. Show that Hk is a cyclic
group under multiplication, and determine its order.
In the following, I denotes the k × k identity matrix.
Let p be an odd prime and let A be a k × k matrix with integer entries. Suppose
that A = I + pm B for some m ∈ N and some k × k matrix B with integer entries, and
that An = I for some n ∈ N. Show that A = I.
Now find the smallest integer r ⩾ 1 such that if A is a k × k matrix with integer
entries satisfying A = I + 2r B for some k × k matrix B with integer entries, and An = I
for some n ∈ N, then in fact A = I.
(i) x2 − 7y 2 = 1.
(ii) x2 − 7y 2 = −1.
(iii) x2 − 7y 2 = 5.
(ii) Consider the following two-step recurrence in n for unm with n ∈ Z+ and
m ∈ Z:
1 h i
un+1
m = (1 − µ)un−1
m + µ(u n
m+1 + u n
m−1 ,
)
1+µ
with µ ⩾ 0. Use Fourier analysis to determine the range of µ for which the
method is stable.
(ii) Find the order of the local truncation error of the scheme.
(iii) In the special case when the matrices B and C commute, does the order of
the scheme change?
N
X
1
IN (h) = h(k/N ).
2N
k=−N +1
∂u ∂ 2 u dw ∂u
= −
∂t ∂x2 dx ∂x
with initial condition u(x, 0) = u0 (x) which is 2-periodic. Seek an approximate solution
for u(x, t) for all t ⩾ 0 that is 2-periodic in x with an expansion of the form
X
u(x, t) = bn (t)eiπnx ,
u
|n|⩽D
bn (t) of the
where D is the truncation level. Write down a differential equation for the u
form
db
un (t) X
= Bnm u bm (t),
dt
|m|⩽D
(i) State, giving a brief justification, necessary and sufficient conditions for xk
to converge to x∗ = A−1 b as k → ∞.
(ii) Give H and v in terms of A and b for the Jacobi method. Prove that the
Jacobi method is convergent if A is strictly diagonally dominant. [You may
use Gershgorin’s circle theorem without proof.]
(i) Show that the unique global minimizer of f is at x∗ = A−1 b. Show also
that for any x we have the identity
(ii) Consider the gradient method with exact line search to find the x that
minimises f (x):
Conclude that
l k
f (xk ) − f (x ) ⩽ 1 −
∗
(f (x0 ) − f (x∗ )) ,
L
(ii) Show that for all k, Ak is symmetric and has the same eigenvalues as A.
A matrix M is r-banded if Mij = 0 whenever |i − j| > r. Show that if Ak
is r-banded then so is Ak+1 .
e k = Q0 . . . Qk−1 and R
(iii) For any k ⩾ 1, let Q ek = Rk−1 . . . R0 . Show that
k e e
A = Qk Rk .
e k . How do these relate to the inverse
(iv) Consider the first and last columns of Q
iteration method and the power method?
with V (t) a real and time-dependent function. Working in the |±⟩ basis, or otherwise,
determine the exact probability to find the system in state |ϕ⟩ at time t > 0 if it was in
state |ψ⟩ at t = 0.
Paper 2, Section II
35B Principles of Quantum Mechanics
(a) State the defining properties of a density operator ρ in a Hilbert space of finite
dimension N , and state the number of real free parameters that determine ρ. Starting
from ρH in the Heisenberg picture, derive the time-dependent ρS (t) in the Schrödinger
picture.
(b) State the commutation relations for the spin operators S with each other and
with S · S . For the pure state |ψ⟩ of a spin- 21 qubit, you are given ⟨Sx ⟩ψ , ⟨Sz ⟩ψ , and the
sign of ⟨Sy ⟩ψ . Determine the normalised state |ψ⟩. [Hint: Recall that the state need only
be determined up to an overall phase, so that the normalised state can be parametrised by
a single complex number.]
(c) For a general mixed state of a spin- 21 qubit, you are given ⟨Sx ⟩, ⟨Sz ⟩, and ⟨Sy ⟩.
Determine the mixed state. Are the expectation values of any three linearly independent
Hermitian operators sufficient to fully specify a general mixed state? Present a proof or a
counterexample.
Hence, or otherwise, prove that if O is a linear operator that commutes with all
components of the angular momentum operator, then
(b) Now consider a two qubit system with states |j1 , m1 ; j2 , m2 ⟩, with j1 = j2 = 1/2.
Let |j, m⟩ be the eigenstates of the total angular momentum of this system. State which
values of j and m occur and write down a resolution of the identity in terms of the |j, m⟩
eigenstates. Write down the values of the Clebsch–Gordan coefficients
j1 , m1 ; j2 , m2 O j1 , m′1 ; j2 , m′2
(i) Compute the corrections to both the ground state of the AB system and
its energy to linear order in λ.
(ii) Using the corrected ground state obtained in part (i), compute the density
operator of the full system and hence the reduced density operator ρA , by
tracing over B. Your final answer should be a density operator ρA that is
normalised up to and including O(λ2 ). [You may assume that corrections
to the ground state beyond linear order do not affect the normalised density
operator up to and including O(λ2 ).]
(iv) Compute the purity of ρA to order λ2 . For what values of λ does it satisfy
the upper and lower bounds? Interpret your findings.
Paper 2, Section II
29L Principles of Statistics
Let X1 , . . . , Xn be i.i.d. observations from a statistical model {f (·, θ) : θ ∈ Θ}
satisfying the usual regularity conditions, where Θ ⊆ Rp . Suppose we wish to test
H0 : θ = θ 0 vs H1 : θ ̸= θ0 .
(a) Define the Wald statistic Wn (θ) and write down a test for H0 based on Wn (θ0 )
with asymptotic type-I error bounded by a given α ∈ (0, 1).
(b) Define the likelihood ratio statistic Λn and write down a test for H0 based on
Λn with asymptotic type-I error bounded by a given α ∈ (0, 1).
(c) Suppose now that p = 1, i.e. θ is a scalar parameter, and we are under the null
H0 . By considering an appropriate Taylor expansion, show that the two test statistics
above are asymptotically equivalent, in the sense that
Λn P
→ 1.
Wn (θ0 )
[You may use, without proof, a uniform law of large numbers, as long as it is clearly
stated.]
(a) Compute the risk R(δMLE , θ), for each θ ∈ (0, 1).
(b) Suppose π is a prior for θ. What is meant by the π-Bayes risk of an estimator
δ of θ? Prove that any estimator of θ which minimises the posterior risk also minimises
the π-Bayes risk.
(c) Now suppose π is the uniform prior on [0, 1]. By differentiating the expression for
the posterior risk with respect to δ(x), prove that δMLE is the unique π-Bayes rule. [In your
calculations, you may interchange R 1 differentiation and integration without justification.
a−1 b−1
You may also use the formulas 0 θ (1 − θ) dθ = Beta(a, b) and Beta(a + 1, b) =
a
Beta(a, b) · a+b without proof.]
(d) Prove that the estimator δMLE is minimax. What does it mean for an estimator
δ of θ to be admissible? Is the estimator δMLE admissible? Justify your answer. [Results
from the course should not be used without proof.]
Tn − θ d
→ N (0, 1).
σn
n−1
vJACK = (An + Bn + 2Cn ),
n
where, for some Rn,i satisfying |Rn,i | ⩽ M (X̄n−1,i − X̄n )2 that you should specify,
2
n
X n
X n
X
An = (g ′ (X̄n ))2 (X̄n−1,i − X̄n )2 , Bn = Rn,i − 1 Rn,j
n
i=1 i=1 j=1
√
and |Cn | ⩽ An Bn . Here X̄n−1,i is the sample mean of the observations with Xi excluded.
P P P
Further assuming that n1 ni=1 (Xi − X̄n )4 → c < ∞, show that vJACK /σn2 → 1.
1
[Hint: Use the fact (which you need not derive) that X̄n−1,i − X̄n = n−1 (X̄n − Xi ).]
(d) Give, with justification, an asymptotically valid (1 − α)-level confidence interval
for θ based on the jackknife estimator.
(a) (i) What does it mean for a measure µ on a measurable space (Ω, F) to be
σ-finite? State the uniqueness of extension theorem for a σ-finite measure.
λ(B) = λ(B + x)
where B + x = {b + x : b ∈ B}.
(iii) Show that Lebesgue measure λ is the unique translation invariant σ-finite
measure on B such that λ((0, 1]) = 1.
(b) Let X, (Xn )n∈N be real-valued random variables with distribution functions FX ,
(FXn )n∈N respectively.
Show that X̃ has the same distribution as X and X̃n has the same
distribution as Xn for all n, and X̃n → X̃ almost surely.
[You may use the fact that for a non-constant, right-continuous, non-
decreasing function g, f (ω) := inf{x ∈ R : ω ⩽ g(x)} is left-continuous
non-decreasing and f (ω) ⩽ x if and only if ω ⩽ g(x). You may also use the
fact that a non-decreasing function has at most a countable set of points of
discontinuity.]
[Hint: You may use without proof the inequality ex ⩾ 1 + x for all x ∈ R.]
(c) Let µ, ν be two measures on a measurable space (Ω, F) such that µ(Ω) < ∞. We
say µ ≪ ν if for any A ∈ F, ν(A) = 0 implies µ(A) = 0. Show that µ ≪ ν if and only if
for all ε > 0 there exists δ > 0 such that for any A ∈ F, µ(A) < ε whenever ν(A) < δ.
Paper 3, Section II
26G Probability and Measure
(iii) Now let (Xn )n∈N be a sequence of real-valued random variables with
characteristic functions (ϕn )n∈N . If for all t ∈ R, ϕn (t) → ϕ(t) for some
function ϕ that is continuous at 0, show that given any ε > 0, there exists
M > 0 such that P(|Xn | > M ) ⩽ ε for all n.
[You may use convergence results for integrals given in the course and Fubini’s
theorem, without proof.]
(a) (i) Write down an expression for the state |Φ1 ⟩ obtained when the state |0⟩⊗n |1⟩
of (n + 1) qubits is acted on by (Uf H ⊗(n+1) ), where H denotes the Hadamard
gate.
(ii) Let |Φ2 ⟩ := H ⊗(n+1) |Φ1 ⟩. Write an expression for this state.
(iii) Let |Φ3 ⟩ := Uf |Φ2 ⟩. Write an expression for this state.
(b) (i) |Φ3 ⟩ is a state of n + 1 qubits. What is the probability of obtaining the n-bit
string a by doing a measurement of the first n qubits in the computational
basis?
(ii) What is the state of the last (i.e., the (n + 1)th ) qubit after the above
measurement?
(iii) Find the probability that a measurement on this qubit yields the value of b
when a contains an odd number of 1s and when a contains an even number
of 1s.
f (x) = 4x mod 9.
Cn , {|βj ⟩}m
j=1 ⊂ C , and a (unique) set of non-negative numbers λ1 , λ2 , . . . , λd (for
m
The number of non-zero λi is called the Schmidt rank of |ψAB ⟩ and is equal to rank(X).
(a) Using the Schmidt decomposition theorem, show that a bipartite state |ψAB ⟩ is
entangled if and only if its Schmidt rank is at least 2.
Here, |ϕ+
BC ⟩ =
√1 (|0B ⟩|0C ⟩
2
+ |1B ⟩|1C ⟩) is a Bell state of two qubits. Show that
(1)
(i) |χABC ⟩ is product across all three bipartitions,
(2)
(ii) |χABC ⟩ is product across the A − BC bipartition, and is entangled across the
B − AC and the C − AB bipartitions,
(3)
(iii) |χABC ⟩ is entangled across all three bipartitions.
(3)
(c) Suppose that Alice (A), Bob (B), and Charlie (C) share either the state |χABC ⟩
(4)
or |χABC ⟩ = √13 (|0A ⟩|0B ⟩|1C ⟩ + |0A ⟩|1B ⟩|0C ⟩ + |1A ⟩|0B ⟩|0C ⟩). Charlie performs a
measurement in the computational basis and obtains the result 0. Show that the
resulting state of Alice and Bob will be
(3)
(i) product in the case they share |χABC ⟩,
(4)
(ii) entangled in the case they share |χABC ⟩.
(I ⊗ A) Φ+ = (AT ⊗ I) Φ+ ,
where |Φ+ ⟩ = √12 (|00⟩ + |11⟩), I denotes the 2 × 2 identity matrix and T denotes
transposition taken with respect to the standard basis {|0⟩ , |1⟩}.
(i) Show that ZX |ψ2 ⟩ = |ψ1 ⟩, where X and Z are the usual one-qubit gates.
(ii) Show that the Bell state |Φ+ ⟩ can be written as √1 (|ψ1 ⟩ |ψ1 ⟩ + |ψ2 ⟩ |ψ2 ⟩).
2
[Hint: Use part (a).]
(iii) Suppose Alice and Bob initially share the Bell state |Φ+ ⟩, with the first
qubit being with Alice and the second qubit being with Bob. Alice then
applies Uθ−1 to her qubit and then measures it in the computational basis.
What are Alice’s possible outcomes and what are the corresponding states
of Bob’s qubit after Alice’s measurement?
(iv) Suppose Alice knows the value of θ but Bob does not. Give a protocol
by which Alice can transmit |ψ1 ⟩ to Bob by sending just one classical bit
to him, given that they share the Bell state |Φ+ ⟩. (This is in contrast to
general teleportation of an unknown qubit, which uses one Bell state and
two classical bits of communication.)
|0⟩ H H
|ψ⟩ U
Here H represents a Hadamard gate, |ψ⟩ denotes a single qubit state which is an eigenstate
of U , U |ψ⟩ = e2πiθ |ψ⟩, and the measurement is done in the computational basis. Express
the phase θ in terms of the probability of the measurement giving zero.
(b) Consider the following circuit which uses a controlled SWAP gate. A SWAP
gate acts as: SWAP |i⟩ |j⟩ 7→ |j⟩ |i⟩ ∀ |i⟩ , |j⟩ ∈ C2 .
|0⟩ H H
|ϕ1 ⟩
SWAP
|ϕ2 ⟩ V
Here V is a unitary operator, |ϕ1 ⟩ , |ϕ2 ⟩ are single qubit states, and the measurement is
done in the computational basis. Find the probability of the measurement giving zero.
(c) Consider the operator U and the state |ψ⟩ introduced in part (a). It is given
that θ = j/2m for some j ∈ {0, 1, 2, . . . , 2m − 1} and some m ∈ N. Define the controlled
unitary operator Θm (U ) which acts on a state |k⟩ |ψ⟩ of m + 1 qubits as:
where U k |ψ⟩ is the state obtained by k successive applications of the operator U on the
state |ψ⟩ and k ∈ {0, 1, 2, . . . , 2m − 1}.
[You should write out the result of applying the operators to |0⟩⊗m |ψ⟩.]
(ii) Write an expression for the corresponding state |ϕj ⟩ of the first m qubits.
m
Show that {|ϕj ⟩}2j=0−1 is an orthonormal basis.
(iv) State the 2 sequential operations that you can do on |ϕj ⟩ to find the value
of j. What is the probability pj of finding the value of j (and hence θ)?
|0⟩ H H
|ψ⟩ U
Here H represents a Hadamard gate, |ψ⟩ denotes a single qubit state which is an eigenstate
of U , U |ψ⟩ = e2πiθ |ψ⟩, and the measurement is done in the computational basis. Express
the phase θ in terms of the probability of the measurement giving zero.
(b) Consider the following circuit which uses a controlled SWAP gate. A SWAP
gate acts as: SWAP |i⟩ |j⟩ 7→ |j⟩ |i⟩ ∀ |i⟩ , |j⟩ ∈ C2 .
|0⟩ H H
|ϕ1 ⟩
SWAP
|ϕ2 ⟩ V
Here V is a unitary operator, |ϕ1 ⟩ , |ϕ2 ⟩ are single qubit states, and the measurement is
done in the computational basis. Find the probability of the measurement giving zero.
(c) Consider the operator U and the state |ψ⟩ introduced in part (a). It is given
that θ = j/2m for some j ∈ {0, 1, 2, . . . , 2m − 1} and some m ∈ N. Define the controlled
unitary operator Θm (U ) which acts on a state |k⟩ |ψ⟩ of m + 1 qubits as:
where U k |ψ⟩ is the state obtained by k successive applications of the operator U on the
state |ψ⟩ and k ∈ {0, 1, 2, . . . , 2m − 1}.
[You should write out the result of applying the operators to |0⟩⊗m |ψ⟩.]
(ii) Write an expression for the corresponding state |ϕj ⟩ of the first m qubits.
m
Show that {|ϕj ⟩}2j=0−1 is an orthonormal basis.
(iv) State the 2 sequential operations that you can do on |ϕj ⟩ to find the value
of j. What is the probability pj of finding the value of j (and hence θ)?
where |ψ1 ⟩ ∈ Sg and |ψ2 ⟩ ∈ Sb . Express a, b, |ψ1 ⟩ and |ψ2 ⟩ in terms of |φ⟩, Pg and Pb .
For any |α⟩ ∈ Hn , consider a unitary operator R|α⟩ which acts as follows:
In particular define R0 := R|0n ⟩ , where |0n ⟩ := |0⟩⊗n . Further, let A denote a unitary
operator which acts on |0n ⟩ as follows:
√ p
|Ω⟩ := A |0n ⟩ := p |ψg ⟩ + 1 − p |ψb ⟩ ,
λ ∈ S(g) =⇒ λk ∈ S(g).
Paper 2, Section II
19H Representation Theory
Let G be a finite group. What is the character χV of a complex representation
(ρ, V ) of G?
Suppose that (ρ, V ) and (σ, W ) are complex representations of G. Show that the
vector space HomC (V, W ) of C-linear maps α : V → W can be made into a representation
of G × G via
((g, h) · α)(v) = σ(h) α(ρ(g −1 )v) for (g, h) ∈ G × G, α ∈ HomC (V, W ) and v ∈ V.
of representations of G × G.
Paper 4, Section II
19H Representation Theory
What is the topological group S 1 ? Assuming any necessary facts about continuous
homomorphisms with domain (R, +), show that every irreducible complex representation
of S 1 is of the form
z 7→ z n : S 1 → GL1 (C)
for some n ∈ Z.
Let ρV : SU (2) → GL(V ) be a complex representation of the topological group
SU (2) and let χV be its character.
(a) Show that χV is determined by its restriction to a subgroup
P T of SU (2)
isomorphic to S 1 and deduce that P
χV may be written in the form n∈Z an z n for some
non-negative integers an such that n∈Z an < ∞. Show moreover that an = a−n for all
n ∈ Z.
(b) Let Vn be the (n + 1)-dimensional irreducible representation of SU (2). Write
down χVn in the form given in part (a). Decompose V4 ⊗ V4 , S 2 V4 and Λ2 V4 as a direct
sum of irreducible representations up to isomorphism.
Paper 2, Section II
24H Riemann Surfaces
For a non-constant analytic map f : R → S between compact Riemann surfaces
and a point z ∈ R, let mf (z) denote the multiplicity of f at z and deg(f ) the degree of f .
State the valency theorem. For the Riemann surface C∞ and a non-constant analytic
function f : C∞ → C∞ , which you may assume is of the form f (z) = p(z)/q(z) for non-
zero polynomials p, q, explain how to find deg(f ). Which f are the analytic isomorphisms
of C∞ ?
If h : C∞ → C∞ is the Möbius transformation that swaps ∞ with 1 and swaps 0
with −1, write down a formula for h, as well as a quadratic equation satisfied by the fixed
points of h.
Now consider the rotational symmetry group G of a regular octahedron P . You may
assume that G is realised as a group of Möbius transformations isomorphic to S4 with the
six vertices of P corresponding to the points 0, ∞, ±1, ±i ∈ C∞ . Write down the possible
sizes of the orbits under this action of G on C∞ .
Consider the function F : C∞ → C∞ given by the formula
Paper 2, Section I
5L Statistical Modelling
Suppose we observe the proportion Y ∼ n−1 Binomial(n, p) where n ∈ N is known
and p ∈ (0, 1) is unknown. Compute the score function and the Fisher information for
this statistical model.
State the Newton–Raphson and Fisher scoring algorithms for computing the max-
imum likelihood estimator.
How many steps do these algorithms take to converge to the maximum likelihood
estimator when initialised at p(0) = Y ? How many steps does Fisher scoring take when
initialised at some p(0) ̸= Y ?
Paper 3, Section I
5L Statistical Modelling
Explain mathematically why the two results returned in the R output below are as
they are.
> n <- 50
> Y <- rnorm(n)
> p <- 2
> Z1 <- matrix(rnorm(n*p), nrow=n)
> lm1 <- lm(Y ~ Z1)
> sum(lm1$residuals)
[1] 0
>
> p <- 49
> Z2 <- matrix(rnorm(n*p), nrow=n)
> lm2 <- lm(Y ~ Z2)
> summary(lm2)$r.squared
[1] 1
Yi = µ + xTi β + εi ,
where µ ∈ R and β ∈ Rp are unknown and the εi are i.i.d. with distribution N (0, σ 2 ) for
known σ 2 > 0. Suppose that it is not possible to measure Yi directly, but we only have
access to a random variable Zi , satisfying
(
1 if Yi > τ,
Zi =
0 otherwise,
y weight width
1 8 3.05 28.3
2 0 1.55 22.5
Coefficients:
Estimate Std. Error t value Pr(>|t|)
(Intercept) -4.5721 4.2132 -1.085 0.2794
weight 1.6817 0.9015 1.865 0.0639
width 0.1446 0.2218 0.652 0.5153
Model 1: y ~ 1
Model 2: y ~ weight + width
Res.Df RSS Df Sum of Sq F Pr(>F)
1 172 1704.9
2 170 1477.7 2 227.2 13.069 5.252e-06 ***
> head(disease)
dis pollution risk
1 0 58.44204 2
2 0 45.29248 1
Coefficients:
Estimate Std. Error z value Pr(>|z|)
(Intercept) -11.83686 3.62821 -3.262 0.00110 **
pollution 0.18243 0.05679 3.212 0.00132 **
risk2 1.53893 0.74665 -2.061 0.03929 *
(a) Write down the generalised linear model being fitted. Why is there no coefficient
for risk1?
(b) Give an interpretation for the coefficient of pollution.
(c) What are the missing degrees of freedom in the output?
(d) How should the R code above be changed to fit the model that corresponds to
the ‘Null deviance’ in the output? What is the AIC value of that model?
(e) State the null and alternative hypotheses for the test performed in the code
below and describe the form of the test.
What can the researcher conclude from the following ouput?
(x2 + y 2 + z 2 )n
U (x) = ,
V 2n/3
where n is a positive integer and V > 0 is an external parameter analogous to volume.
(i) Calculate the partition function and hence show that the Helmholtz free
energy is
h i
F = −N kB T ln V + A ln(kB T ) + ln In + B ,
where Z ∞
2n
In = u2 e−u du ,
0
and you should determine A and B.
∂F
(ii) Considering the conjugate pressure to V , p = − ∂V T,N , derive the
equation of state.
(iii) Compute the average energy E, the variance of energy (∆E)2 and the heat
capacity CV for the system. Comment on the behaviour of (∆E)/E in the
thermodynamic limit.
(iv) Obtain the local particle number density as a function of x and hence
determine the most likely |x| to find a particle.
(c) Explain why the energy E(S, V, N ) is a homogeneous function of degree 1, where
S is the entropy, V is the volume and N is the number of particles. Hence, using the
first law of thermodynamics, find an expression for E in terms of S, V , N , µ, p and T ,
where µ is the chemical potential, p is the pressure and T is the temperature. Show that
dµ = (V dp − SdT )/N .
(d) Consider a chemical reaction at constant T and p where each molecule of
chemical A can change into two molecules of chemical B and one molecule of chemical C,
and vice-versa, i.e. A ↔ 2B + C. By minimising the Gibbs free energy G, derive a relation
between the chemical potentials of the three chemicals at equilibrium, where the chemical
∂G
potential of chemical i is µi = ∂N i
.
(iii) Find the critical temperature Tc for the Dieterici equation of state in terms
of a, b and kB .
and ( )
n
X
Q= q ∈ Rn : P ⊤ q = 0, qi = 1, qi > 0 for all i .
i=1
Consider a one-period market model with d risky assets, where Sti is the price of asset
i at time t ∈ {0, 1} and r is the interest rate. Assume that there exists at least one risk-
neutral measure for the model, and that the random variables {S1i − (1 + r)S0i : 1 ⩽ i ⩽ d}
are linearly independent.
(c) Let Y be a random variable such that EQ (Y ) > 0 for all risk-neutral measures
Q. Show that there exists a vector θ ∈ Rd such that
For any 0 < z < 1, compute E(z T ). [You may use the fact that T is a finite stopping time
without proof.]
P(ξ1 = 1 + b) = p = 1 − P(ξ1 = 1 + a)
for constants −1 < a < b and 0 < p < 1. Assume that there exists a risk-neutral measure
Q for the model.
(a) Show that a < r < b, and that the sequence (ξn )n⩾1 is IID under Q. Find
Q(ξ1 = 1 + b).
Fix a maturity date N > 0 and a payout function g. For all s > 0 define
V (N, s) = g(s),
q 1−q
V (n − 1, s) = V (n, s(1 + b)) + V (n, s(1 + a)) for all 1 ⩽ n ⩽ N,
1+r 1+r
where q = Q(ξ1 = 1 + b).
(b) Consider the case where g(s) = log s. Find, with justification, a formula for
V (n, s) of the form V (n, s) = An log s + Bn , for families of constants (An )0⩽n⩽N and
(Bn )0⩽n⩽N which should be specified.
(c) Prove that an investor with initial capital X0 = V (0, S0 ) can replicate the payout
of the vanilla European claim with time-N payout g(SN ). How many shares of the stock
should the investor hold during the time interval (n − 1, n] ?
Now assume a ⩽ 0, and fix a barrier B > S0 . For all s > 0 define
U (N, s) = g(s),
q 1−q
U (n − 1, s) = U (n, s(1 + b))1{s(1+b)<B} + U (n, s(1 + a)) for all 1 ⩽ n ⩽ N.
1+r 1+r
(d) In terms of the function U , derive an explicit formula for the number of shares of
the stock an investor should hold during the time interval (n−1, n] in order to replicate the
payout of the up-and-out European claim with time-N payout g(SN )1{max0⩽n⩽N Sn <B} .
Consider a continuous time market with interest rate r and a stock whose time-t
price is
1 2
St = S0 e(r− 2 σ )t+σWt ,
for a given volatility parameter σ > 0.
(c) Show that the discounted stock price (e−rt St )t⩾0 is a martingale with respect to
the filtration generated by the Brownian motion.
(d) Show that the time-0 Black–Scholes price of a European call with strike K and
maturity date T is S0 F (σ 2 T, Ke−rT /S0 ), where
1 √
F (v, m) = E[(e− 2 v+ vZ
− m)+ ] for all v, m ⩾ 0
1 2 )T 1 2 )T
S0 e−( 2 r+ασ F ( 13 σ 2 T, Ke−( 2 r−ασ /S0 ),
for a constant α to be determined. [You may assume that the process (It )t⩾0 defined in
part (b) is Gaussian.]
Paper 2, Section I
2G Topics in Analysis
State Baire’s category theorem. Define an isolated point for a metric space.
Which of the following statements are true and which are false? Give a proof or a
counterexample. (By a metric space we mean a non-empty metric space.)
(vi) All the points in a countable complete metric space are isolated.
(vii) A countable complete metric space containing at least two points must
contain at least two isolated points.
(viii) A countable complete metric space containing infinitely many points must
contain infinitely many isolated points.
Paper 4, Section I
2G Topics in Analysis
(a) Show that the collection of algebraic numbers is countable.
P
(b) Suppose that aj > 0 and ∞ j=1 aj converges. Show that we can find θj ∈ {0, 1}
so that
X∞
θj aj
j=1
is transcendental.
(c) Suppose θj ∈ {0, 1} and θj = 1 for infinitely many values of j. Show that
∞
X θj
j!
j=1
is irrational.
pn b0
= a0 +
qn b1
a1 +
b2
a2 +
b3
a3 +
..
.
bn−1
an−1 +
an
and that
pn bn pn−1 a0 b0 a1 b1 an bn
= ... .
qn bn qn−1 1 0 1 0 1 0
We now specialise to the case when bj = 1 and the aj are strictly positive integers
for all j. Show that pn qn−1 − qn pn−1 = (−1)n+1 for all n ⩾ 1.
Show that pn /qn tends to a limit x and
pn pn+1 1
−x + −x =
qn qn+1 qn qn+1
for each n ⩾ 0.
Now specialise still further to the case when aj = 1 for all j. Show that pn = Fn+2 ,
qn = Fn+1 for all n ⩾ 1 where F0 = 0, F1 = 1 and Fn+2 = Fn+1 + Fn . Solve√ this difference
equation to obtain an expression for Fn in terms of powers of ϕ = (1 + 5)/2. Hence show
that, in this case, the limit x discussed in the previous paragraph is ϕ. Show also that
Fn+1 ϕ Fn+2 1
Fn Fn+1 −ϕ → √ and Fn Fn+1 −ϕ → √
Fn 5 Fn+1 ϕ 5
as n → ∞.
(i) Suppose that each vertex of the grid is coloured red, green or blue, that every
vertex on DE is coloured red or green, every vertex on EF is coloured green
or blue and every vertex on F D is coloured blue or red. Show that there is
an odd number of triangles in the grid with vertices of different colours.
(ii) Suppose instead that each vertex of the grid is coloured red, green or blue,
and that the three vertices D, E and F have different colours. Must there
be a triangle in the grid with vertices of different colours? Give reasons.
△ = {λ1 d + λ2 e + λ3 f : λ1 + λ2 + λ3 = 1, λu ⩾ 0},
where d, e, f are the position vectors of D, E and F , and write I for the closed line
segment DE, J for the closed line segment EF , K for the closed line segment F D.
Which of the following statements are true and which are false? Give a proof or a
counterexample.
(i) Show that the complex amplitude of the velocity potential in x > 0 is given
by
i ρρ−
+
sin λ − cc−
+
cos λ ωL
ϵc− 2 2 , where λ= .
ρ+ 2 c− 2λ c−
ρ− sin λ + c+ cos
(ii) Consider the two sets of frequencies of oscillation such that λ = nπ and
λ = (n + 12 )π for integer n. Calculate the time-averaged acoustic energy
flux in x > 0 for each set, and briefly comment on the behaviour in the case
where ρ+ ≪ ρ− and c+ ≈ c− .
∂2u
ρ = (λ + µ)∇(∇ · u) + µ∇2 u,
∂t2
where the constants λ and µ are the Lamé moduli.
(a) Show that this equation supports two distinct classes of wave-like motion: P-
waves for the dilatation ϑ = ∇ · u with phase speed cp ; and S-waves for the rotation
Ω = ∇ × u with phase speed cs . You should express cp and cs explicitly in terms of the
Lamé moduli.
(b) Now consider a region of this solid with a horizontal plane boundary at z = 0
in which plane waves propagate with wave vector k = κ(sin θ, 0, cos θ), i.e. θ is the angle
the wave vector makes with the vertical z−direction and κ is the magnitude of the wave
vector. Explain briefly why such a domain can in general support:
You should define explicitly the orientations of the complex vector amplitudes A, BV and
BH .
(c) Now consider a region of the solid between a rigid plane boundary at z = 0 and
a free surface at z = h > 0.
(i) Show that this region can support propagating SH-waves with wave vector
in the x−direction (i.e. with θ = π/2), calculating explicitly the dispersion
relation.
(ii) Deduce that there is a cut-off frequency ωn for each mode in the vertical,
given by
(2n + 1)π
ωn = cs ,
2h
for non-negative integers n = 0, 1, 2 . . . .
(iii) Express the phase velocity c and the group velocity cg of each mode in terms
of cs , ωn and κ.
(iv) Deduce that, for any given wave number κ > 0 and mode with n ⩾ 0,
c = mcg with m > 1, where you should express m in terms of n, κ and h.
(v) Calculate m explicitly for the specific wave with horizontal wavelength h
and n = 1.
1 dρ0 ∂w ∂2w
≪ .
ρ0 dz ∂z ∂z 2
(ii) You may assume that the vertical component of the fluid velocity at z = 0 is
given by the vertical component of the velocity of fluid particles that follow
the undulations of the lower boundary, i.e. w = U ∂h/∂x at z = 0. Hence
establish that for z > 0, the vertical velocity satisfies
w = w0 exp[i(kx + mz − ωt)],
(c − c0 )
R± = u ± 2 ,
(γ − 1)
dx
= u ± c,
dt
p
where u is the fluid velocity, and c = dp/dρ is the associated sound speed.
(b) Consider a semi-infinite tube filled with such an ideal gas in the region x > X(t)
to the right of a piston at x = X(t). At time t = 0, the piston and the gas are at rest,
X = 0, and the gas is uniform with c = c0 . For t > 0, the piston accelerates smoothly in
the positive x−direction.
(i) Show that, prior to the formation of a shock, the motion of the gas can be
written in terms of a parameter τ by
1
u(x, t) = Ẋ(τ ) on x = X(τ ) + c0 + (γ + 1)Ẋ(τ ) (t − τ ),
2
c0 t3
X(t) = ,
T2
where T is a positive constant. Show that a shock first forms in the gas
when t/T = f (γ), for some function f (γ) which you should determine.