Physics Basic
Physics Basic
Boundless
Boundless
Physics
Boundless
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TABLE OF CONTENTS
Licensing
2: Kinematics
2.1: Basics of Kinematics
2.2: Speed and Velocity
2.3: Acceleration
2.4: Problem-Solving for Basic Kinematics
2.5: Free-Falling Objects
3: Two-Dimensional Kinematics
3.1: Motion in Two Dimensions
3.2: Vectors
3.3: Projectile Motion
3.4: Multiple Velocities
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6.4: Work-Energy Theorem
6.5: Potential Energy and Conservation of Energy
6.6: Power
6.7: CASE STUDY: World Energy Use
6.8: Further Topics
10: Fluids
10.1: Introduction
10.2: Density and Pressure
10.3: Archimedes’ Principle
10.4: Cohesion and Adhesion
10.5: Fluids in Motion
10.6: Deformation of Solids
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12: Temperature and Kinetic Theory
12.1: Introduction
12.2: Temperature and Temperature Scales
12.3: Thermal Expansion
12.4: Ideal Gas Law
12.5: Kinetic Theory
12.6: Phase Changes
12.7: The Zeroth Law of Thermodynamics
12.8: Thermal Stresses
12.9: Diffusion
14: Thermodynamics
14.1: Introduction
14.2: The First Law of Thermodynamics
14.3: The Second Law of Thermodynamics
14.4: Entropy
14.5: The Third Law of Thermodynamics
16: Sound
16.1: Introduction
16.2: Sound Intensity and Level
16.3: Doppler Effect and Sonic Booms
16.4: Interactions with Sound Waves
16.5: Further Topics
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18: Electric Potential and Electric Field
18.1: Overview
18.2: Equipotential Surfaces and Lines
18.3: Point Charge
18.4: Capacitors and Dielectrics
18.5: Applications
21: Magnetism
21.1: Magnetism and Magnetic Fields
21.2: Magnets
21.3: Magnetic Force on a Moving Electric Charge
21.4: Motion of a Charged Particle in a Magnetic Field
21.5: Magnetic Fields, Magnetic Forces, and Conductors
21.6: Applications of Magnetism
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25: Vision and Optical Instruments
25.1: The Human Eye
25.2: Other Optical Instruments
Index
Glossary
Detailed Licensing
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Licensing
A detailed breakdown of this resource's licensing can be found in Back Matter/Detailed Licensing.
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CHAPTER OVERVIEW
This page titled 1: The Basics of Physics is shared under a not declared license and was authored, remixed, and/or curated by Boundless.
1
1.1: The Basics of Physics
Introduction: Physics and Matter
Physics is a study of how the universe behaves.
learning objectives
Apply physics to describe the function of daily life
Physics is a natural science that involves the study of matter and its motion through space and time, along with related concepts
such as energy and force. More broadly, it is the study of nature in an attempt to understand how the universe behaves.
What is Physics?
What is Physics?: Mr. Andersen explains the importance of physics as a science. History and virtual examples are used to give the
discipline context.
Physics uses the scientific method to help uncover the basic principles governing light and matter, and to discover the implications
of those laws. It assumes that there are rules by which the universe functions, and that those laws can be at least partially
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understood by humans. It is also commonly believed that those laws could be used to predict everything about the universe’s future
if complete information was available about the present state of all light and matter.
Matter is generally considered to be anything that has mass and volume. Many concepts integral to the study of classical physics
involve theories and laws that explain matter and its motion. The law of conservation of mass, for example, states that mass cannot
be created or destroyed. Further experiments and calculations in physics, therefore, take this law into account when formulating
hypotheses to try to explain natural phenomena.
Physics aims to describe the function of everything around us, from the movement of tiny charged particles to the motion of
people, cars, and spaceships. In fact, almost everything around you can be described quite accurately by the laws of physics.
Consider a smart phone; physics describes how electricity interacts with the various circuits inside the device. This knowledge
helps engineers select the appropriate materials and circuit layout when building the smart phone. Next, consider a GPS system;
physics describes the relationship between the speed of an object, the distance over which it travels, and the time it takes to travel
that distance. When you use a GPS device in a vehicle, it utilizes these physics equations to determine the travel time from one
location to another. The study of physics is capable of making significant contributions through advances in new technologies that
arise from theoretical breakthroughs.
Global Positioning System: GPS calculates the speed of an object, the distance over which it travels, and the time it takes to travel
that distance using equations based on the laws of physics.
learning objectives
Explain why the study of physics is integral to the study of other sciences
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Physics in Chemistry: The study of matter and electricity in physics is fundamental towards the understanding of concepts in
chemistry, such as the covalent bond.
Physics has many applications in the biological sciences. On the microscopic level, it helps describe the properties of cell walls and
cell membranes. On the macroscopic level, it can explain the heat, work, and power associated with the human body. Physics is
involved in medical diagnostics, such as X-rays, magnetic resonance imaging (MRI), and ultrasonic blood flow measurements.
Medical therapy sometimes directly involves physics: cancer radiotherapy uses ionizing radiation, for instance. Physics can also
explain sensory phenomena, such as how musical instruments make sound, how the eye detects color, and how lasers can transmit
information.
The boundary between physics and the other sciences is not always clear. For instance, chemists study atoms and molecules, which
are what matter is built from, and there are some scientists who would be equally willing to call themselves physical chemists or
chemical physicists. It might seem that the distinction between physics and biology would be clearer, since physics seems to deal
with inanimate objects. In fact, almost all physicists would agree that the basic laws of physics that apply to molecules in a test tube
work equally well for the combination of molecules that constitutes a bacterium. What differentiates physics from biology is that
many of the scientific theories that describe living things ultimately result from the fundamental laws of physics, but cannot be
rigorously derived from physical principles.
It is not necessary to formally study all applications of physics. What is most useful is the knowledge of the basic laws of physics
and skill in the analytical methods for applying them. The study of physics can also improve your problem-solving skills.
Furthermore, physics has retained the most basic aspects of science, so it is used by all of the sciences. The study of physics makes
other sciences easier to understand.
learning objectives
Define the terms model, theory, and law
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Models
A model is a representation of something that is often too difficult (or impossible) to display directly. While a model’s design is
justified using experimental information, it is only accurate under limited situations. An example is the commonly used “planetary
model” of the atom, in which electrons are pictured as orbiting the nucleus, analogous to the way planets orbit the Sun. We cannot
observe electron orbits directly, but the mental image helps explain the observations we can make, such as the emission of light
from hot gases. Physicists use models for a variety of purposes. For example, models can help physicists analyze a scenario and
perform a calculation, or they can be used to represent a situation in the form of a computer simulation.
Planetary Model of an Atom: The planetary model of the atom in which electrons are pictured as orbiting the nucleus, analogous
to the way planets orbit the Sun
Theories
A theory is an explanation for patterns in nature that is supported by scientific evidence and verified multiple times by various
groups of researchers. Some theories include models to help visualize phenomena, whereas others do not. Newton’s theory of
gravity, for example, does not require a model or mental image, because we can observe the objects directly with our own senses.
The kinetic theory of gases, on the other hand, makes use of a model in which a gas is viewed as being composed of atoms and
molecules. Atoms and molecules are too small to be observed directly with our senses—thus, we picture them mentally to
understand what our instruments tell us about the behavior of gases.
Laws
A law uses concise language to describe a generalized pattern in nature that is supported by scientific evidence and repeated
experiments. Often, a law can be expressed in the form of a single mathematical equation. Laws and theories are similar in that
they are both scientific statements that result from a tested hypothesis and are supported by scientific evidence. However, the
designation law is reserved for a concise and very general statement that describes phenomena in nature, such as the law that
energy is conserved during any process, or Newton’s second law of motion, which relates force, mass, and acceleration by the
simple equation F = ma . A theory, in contrast, is a less concise statement of observed phenomena. For example, the Theory of
Evolution and the Theory of Relativity cannot be expressed concisely enough to be considered a law. The biggest difference
between a law and a theory is that a law is much more complex and dynamic, and a theory is more explanatory. A law describes a
single observable point of fact, whereas a theory explains an entire group of related phenomena. And, whereas a law is a postulate
that forms the foundation of the scientific method, a theory is the end result of that process.
Key Points
Physics is a natural science that involves the study of matter and its motion through space and time, along with related concepts
such as energy and force.
Matter is generally considered to be anything that has mass and volume.
Scientific laws and theories express the general truths of nature and the body of knowledge they encompass. These laws of
nature are rules that all natural processes appear to follow.
Many scientific disciplines, such as biophysics, are hybrids of physics and other sciences.
The study of physics encompasses all forms of matter and its motion in space and time.
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The application of physics is fundamental towards significant contributions in new technologies that arise from theoretical
breakthroughs.
Concepts in physics cannot be proven, they can only be supported or disproven through observation and experimentation.
A model is an evidence-based representation of something that is either too difficult or impossible to display directly.
A theory is an explanation for patterns in nature that is supported by scientific evidence and verified multiple times by various
groups of researchers.
A law uses concise language, often expressed as a mathematical equation, to describe a generalized pattern in nature that is
supported by scientific evidence and repeated experiments.
Key Terms
matter: The basic structural component of the universe. Matter usually has mass and volume.
scientific method: A method of discovering knowledge about the natural world based in making falsifiable predictions
(hypotheses), testing them empirically, and developing peer-reviewed theories that best explain the known data.
application: the act of putting something into operation
Model: A representation of something difficult or impossible to display directly
Law: A concise description, usually in the form of a mathematical equation, used to describe a pattern in nature
theory: An explanation for patterns in nature that is supported by scientific evidence and verified multiple times by various
groups of researchers
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1.2: Units
Length
Length is a physical measurement of distance that is fundamentally measured in the SI unit of a meter.
learning objectives
Distinguish SI and customary units of length
Length can be defined as a measurement of the physical quantity of distance. Many qualitative observations fundamental to physics
are commonly described using the measurement of length. The distance between objects, the rate at which objects are traveling,
and how much force an object exerts are all dependent on length as a variable. In order to describe length in a standardized and
quantitative manner, an accepted unit of measurement must be utilized.
Many different units of length are used around the world. In the United States, the U.S. customary units operationally describe
length in terms of the basic unit of an inch. Varying lengths are thus described in relation to the inch, such as a foot equaling 12
inches, a yard equaling three feet, and a mile equaling 1,760 yards.
Though regional use of different measurement units is not generally problematic, it can raise issues of compatibility and
understanding when working abroad or collaboratively with international partners. As such, a standard unit of measurement that is
internationally accepted is needed. The basic unit of length as identified by the International System of Units (SI) is the meter. The
meter is expressed more specifically in terms of speed of light.
One meter is defined as the distance that light travels in a vacuum in 1
299,792,458
of a second. All lengths are measured in terms
related to the meter, where its multiples are devised around the convenience of the number 10. For example, a centimeter is equal to
1
100
of a meter (or 10 meters), and a kilometer is equal to 1,000 meters (or 10 meters).
−2 3
Meter Defined by Speed of Light: The meter is defined to be the distance that light travels in 1
299,792,458
of a second in a vacuum.
Distance traveled is speed multiplied by time.
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High School Physics - Metric System
Metric System – Length: A brief introduction to the metric system and unit conversions.
Mass
Mass is the quantity of matter that an object contains, as measured by its resistance to acceleration.
learning objectives
Explain the difference between mass and weight
Mass
Mass, specifically inertial mass, is a quantitative measure of an object’s resistance to acceleration. It is an intrinsic property of an
object and does not change because of the environment. The SI unit of mass is the kilogram (kg).
The kilogram is defined as being equal to the mass of the International Prototype Kilogram (IPK), which is almost exactly equal to
the mass of one liter of water. It is also the only SI unit that is directly defined by an artifact, rather than a fundamental physical
property that can be reproduced in different laboratories. Four of the seven base units in the SI system are defined relative to the
kilogram, so the stability of this measurement is crucial for accurate and consistent measurements.
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In 2005, the International Committee for Weights and Measures (CIPM) recommended that the kilogram be redefined in terms of a
fundamental constant of nature, due to evidence that the International Prototype Kilogram will vary in mass over time. At its 2011
meeting, the General Conference on Weights and Measures (CGPM) agreed that the kilogram should be redefined in terms of the
Planck constant. The conference deferred a final decision until its next meeting in 2014.
Prototype Mass Drifts: A graph of the relative change in mass of selected kilogram prototypes.
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High School Physics - Metric System
Metric System – Mass: A brief introduction to the metric system and unit conversions.
Time
Time is the fundamental physical quantity of duration and is measured by the SI Unit known as the second.
learning objectives
Relate time with other physical quantities
Time is one of the seven fundamental physical quantities in the International System (SI) of Units. Time is used to define other
quantities, such as velocity or acceleration, and as such, it is important that it be standardized and quantified precisely. An
operational definition of time is highly useful in the conduct of both advanced experiments and everyday affairs of life.
Historically, temporal measurement was a prime motivation in navigation and astronomy. Periodic events and motion have long
served as standards for units of time. For example, the movement of the sun across the sky, the phases of the moon, the swing of a
pendulum, and the beat of a heart have all been used as a standard for time keeping. These events and standards, however, are
highly dynamic in nature and cannot reliably be utilized for accurate quantitative measures. Between 1000 and 1960 the second
was defined as 1
86,400
of a mean solar day. This definition changed between 1960 and 1967 and was defined in terms of the period
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of the Earth’s orbit around the Sun in 1900. Today, the SI Unit of the second is defined in terms of radiation emitted by cesium
atoms.
The second is now operationally defined as “the duration of 9,192,631,770 periods of the radiation corresponding to the transition
between the two hyperfine levels of the ground state of the cesium 133 atom.” It follows that the hyperfine splitting in the ground
state of the cesium 133 atom is exactly 9,192,631,770 hertz. In other words, cesium atoms can be made to vibrate in a very steady
way, and these vibrations can be readily observed and counted. The second is the time required for 9,192,631,770 of these
vibrations to occur.
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NIST-F1 Cesium Clock: NIST-F1 is referred to as a fountain clock because it uses a fountain-like movement of atoms to obtain its
improved reckoning of time.
Metric System – Time: A brief introduction to the metric system and unit conversions.
learning objectives
Apply prefixes to units and distinguish between SI and customary units
Prefixes
A metric prefix, or SI prefix, is a unit prefix that precedes a basic unit of measure to indicate a multiple or fraction of the unit. Each
prefix has a unique symbol that is prepended to the unit symbol. The prefix kilo-, for example, may be added to gram to indicate
multiplication by one thousand; one kilogram is equal to one thousand grams (1 kg = 1000 g). The prefix centi-, likewise, may be
added to meter to indicate division by one hundred; one centimeter is equal to one hundredth of a meter (1 cm = 0.01 m). Prefixes
in varying multiples of 10 are a feature of all forms of the metric system, with many dating back to the system’s introduction in the
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1790s. Today, the prefixes are standardized for use in the International System of Units (SI) by the International Bureau of Weights
and Measures.There are twenty prefixes officially specified by SI.
SI Unit Prefixes: The twenty prefixes officially specified by the International System of Units
It is important to note that the kilogram is the only SI unit with a prefix as part of its name and symbol. Because multiple prefixes
may not be used, in the case of the kilogram the prefix names are used with the unit name “gram” and the prefix symbols are used
with the unit symbol “g. ” With this exception, any SI prefix may be used with any SI unit, including the degree Celsius and its
symbol °C.
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High School Physics - Metric System
Metric System – Prefixes: A brief introduction to the metric system and unit conversions.
Converting Units
Converting between units can be done through the use of conversion factors or specific conversion formulas.
learning objectives
Apply factor-label method for converting units
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Unit Conversion in the Metric System - CLEAR & SI…
SI…
Unit Conversion in the Metric System: EASY Unit Conversion in the Metric System – This simple extra help video tutorial
explains the metric system and how to make simple metric conversions.
Conversion Methods
There are several ways to approach doing conversions. One commonly used method is known as the Factor-label method for
converting units, or the “railroad method. ”
The factor-label method is the sequential application of conversion factors expressed as fractions and arranged so that any
dimensional unit appearing in both the numerator and denominator of any of the fractions can be cancelled out until only the
desired set of dimensional units is obtained. For example, 10 miles per hour can be converted to meters per second by using a
sequence of conversion factors.
Each conversion factor is equivalent to the value of one. For example, starting with 1 mile = 1609 meters and dividing both sides of
the equation by 1 mile yields 1 mile
1 mile
=
1609 meters
1 mile
, which when simplified yields 1 = 1609 meters
1 mile
. Physically crossing out the units
that cancel each other out will also help visualize what’s left over.
Converting 1 year into seconds using the Factor-Label Method: Physically crossing out units that cancel out helps visualize the
“leftover” unit(s).
So, when the units mile and hour are cancelled out and the arithmetic is done, 10 miles per hour converts to 4.47 meters per second.
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A limitation of the factor-label method is that it can only convert between units that have a constant ratio that can be multiplied, or
a multiplication factor. This method cannot be used between units that have a displacement, or difference factor. An example is the
conversion between degrees Celsius and kelvins, or between Celsius and Fahrenheit. For these, it is best to use the specific
conversion formulas.
For example, if you are planning a trip abroad in Spain and the weather forecast predicts the weather to be mostly cloudy and 16°C,
you may want to convert the temperature into °F, a unit that you are more comfortable interpreting. In order to do this, you would
need to know the conversion formula from Celsius to Fahrenheit. This formula is: [°F] = [°C] × + 32 .
9
9
[°F] = [°C] × + 32 (1.2.1)
5
So you would then know that 16°C is equivalent to 60.8°F and be able to pack the right type of clothing to be comfortable.
Key Points
The SI unit for length is the meter.
One meter is defined as the distance that light travels in a vacuum in 1
of a second.
299,792,458
Derivatives of measurement units related to the meter are devised around the convenience of the number 10.
The kilogram is the only SI unit directly defined by the artifact itself.
Mass is a property that does not depend on gravitational fields, unlike weight.
One kilogram is defined as the mass of the International Prototype Kilogram (IPK), a platinum-iridium alloy cylinder.
One kilogram is almost exactly equal to the mass of one liter of water.
Time is a physical quantity of duration.
The SI Unit for time is the second.
The second is operationally defined in terms of radiation emitted by cesium atoms.
The twenty standardized prefixes for use in the International System of Units are derived from multiples of 10.
The kilogram is the only SI unit with a prefix as part of its name and symbol; as such, SI unit prefixes are prepended to the unit
gram.
The United States customary units define measurements based on the English, or Imperial, unit standards.
Conversion of units is the conversion between different units of measurement for the same quantity, typically through
multiplicative conversion factors.
The factor-label method is the sequential application of conversion factors expressed as fractions in which units appearing in
both the numerator and denominator can be cancelled out, leaving only the desired set of units.
For conversions that have a difference factor, specific conversion formulas should be used.
Key Terms
Length: How far apart objects are physically.
acceleration: the rate at which the velocity of a body changes with time
inertia: the tendency of an object to resist any change in its motion
Radiation: the emission of energy as electromagnetic waves or as moving or oscillating subatomic particles.
prefix: That which is prefixed; especially one or more letters or syllables added to the beginning of a word to modify its
meaning; as, pre- in prefix, con- in conjure.
conversion: a change between different units of measurement for the same quantity.
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1.3: Significant Figures and Order of Magnitude
Scientific Notation
Scientific notation is a way of writing numbers that are too big or too small in a convenient and standard form.
learning objectives
Convert properly between standard and scientific notation and identify appropriate situations to use it
Scientific Notation: There are three parts to writing a number in scientific notation: the coefficient, the base, and the exponent.
Most of the interesting phenomena in our universe are not on the human scale. It would take about 1,000,000,000,000,000,000,000
bacteria to equal the mass of a human body. Thomas Young’s discovery that light was a wave preceded the use of scientific
notation, and he was obliged to write that the time required for one vibration of the wave was “ of a millionth of a millionth of a
500
1
second”; an inconvenient way of expressing the point. Scientific notation is a less awkward and wordy way to write very large and
very small numbers such as these.
A Simple System
Scientific notation means writing a number in terms of a product of something from 1 to 10 and something else that is a power of
ten.
For instance, 32 = 3.2 ⋅ 10 1
2
320 = 3.2 ⋅ 10
3200 = 3.2 ⋅ 10
3
, and so forth…
Each number is ten times bigger than the previous one. Since 10 is ten times smaller than 10 , it makes sense to use the notation
1 2
10 to stand for one, the number that is in turn ten times smaller than 10 . Continuing on, we can write 10 to stand for 0.1, the
0 1 −1
number ten times smaller than 10 . Negative exponents are used for small numbers:
0
0
3.2 = 3.2 ⋅ 10
−1
0.32 = 3.2 ⋅ 10
−2
0.032 = 3.2 ⋅ 10
Scientific notation displayed calculators can take other shortened forms that mean the same thing. For example, 3.2 ⋅ 10
6
(written
notation) is the same as 3.2E + 6 (notation on some calculators) and 3.2 (notation on some other calculators).
6
1.3.1 https://phys.libretexts.org/@go/page/14433
Round-off Error
A round-off error is the difference between the calculated approximation of a number and its exact mathematical value.
learning objectives
Explain the impact round-off errors may have on calculations, and how to reduce this impact
Round-off Error
A round-off error, also called a rounding error, is the difference between the calculated approximation of a number and its exact
mathematical value. Numerical analysis specifically tries to estimate this error when using approximation equations, algorithms, or
both, especially when using finitely many digits to represent real numbers. When a sequence of calculations subject to rounding
errors is made, errors may accumulate, sometimes dominating the calculation.
Calculations rarely lead to whole numbers. As such, values are expressed in the form of a decimal with infinite digits. The more
digits that are used, the more accurate the calculations will be upon completion. Using a slew of digits in multiple calculations,
however, is often unfeasible if calculating by hand and can lead to much more human error when keeping track of so many digits.
To make calculations much easier, the results are often ’rounded off’ to the nearest few decimal places.
For example, the equation for finding the area of a circle is A = πr . The number π(pi) has infinitely many digits, but can be
2
truncated to a rounded representation of as 3.14159265359. However, for the convenience of performing calculations by hand, this
number is typically rounded even further, to the nearest two decimal places, giving just 3.14. Though this technically decreases the
accuracy of the calculations, the value derived is typically ‘close enough’ for most estimation purposes.
However, when doing a series of calculations, numbers are rounded off at each subsequent step. This leads to an accumulation of
errors, and if profound enough, can misrepresent calculated values and lead to miscalculations and mistakes.
The following is an example of round-off error:
−−−−−−−−− − −−−− −−−−−
2 2
√4.58 + 3.28 = √21.0 + 10.8 = 5.64
Rounding these numbers off to one decimal place or to the nearest whole number would change the answer to 5.7 and 6,
respectively. The more rounding off that is done, the more errors are introduced.
learning objectives
Choose when it is appropriate to perform an order-of-magnitude calculation
Orders of Magnitude
An order of magnitude is the class of scale of any amount in which each class contains values of a fixed ratio to the class preceding
it. In its most common usage, the amount scaled is 10, and the scale is the exponent applied to this amount (therefore, to be an
order of magnitude greater is to be 10 times, or 10 to the power of 1, greater). Such differences in order of magnitude can be
measured on the logarithmic scale in “decades,” or factors of ten. It is common among scientists and technologists to say that a
parameter whose value is not accurately known or is known only within a range is “on the order of” some value. The order of
magnitude of a physical quantity is its magnitude in powers of ten when the physical quantity is expressed in powers of ten with
one digit to the left of the decimal.
Orders of magnitude are generally used to make very approximate comparisons and reflect very large differences. If two numbers
differ by one order of magnitude, one is about ten times larger than the other. If they differ by two orders of magnitude, they differ
by a factor of about 100. Two numbers of the same order of magnitude have roughly the same scale — the larger value is less than
ten times the smaller value.
It is important in the field of science that estimates be at least in the right ballpark. In many situations, it is often sufficient for an
estimate to be within an order of magnitude of the value in question. Although making order-of-magnitude estimates seems simple
and natural to experienced scientists, it may be completely unfamiliar to the less experienced.
1.3.2 https://phys.libretexts.org/@go/page/14433
Example 1.3.1:
Some of the mental steps of estimating in orders of magnitude are illustrated in answering the following example question:
Roughly what percentage of the price of a tomato comes from the cost of transporting it in a truck?
Guessing the Number of Jelly Beans: Can you guess how many jelly beans are in the jar? If you try to guess directly, you
will almost certainly underestimate. The right way to do it is to estimate the linear dimensions and then estimate the volume
indirectly.
Incorrect solution: Let’s say the trucker needs to make a prot on the trip. Taking into account her benits, the cost of gas, and
maintenance and payments on the truck, let’s say the total cost is more like 2000. You might guess about 5000 tomatoes would
t in the back of the truck, so the extra cost per tomato is 40 cents. That means the cost of transporting one tomato is
comparable to the cost of the tomato itself.
The problem here is that the human brain is not very good at estimating area or volume — it turns out the estimate of 5000
tomatoes fitting in the truck is way off. (This is why people have a hard time in volume-estimation contests, such as the one
shown below.) When estimating area or volume, you are much better off estimating linear dimensions and computing the
volume from there.
So, here’s a better solution: As before, let’s say the cost of the trip is $2000. The dimensions of the bin are probably 4m by 2m
by 1m, for a volume of 8 m . Since our goal is just an order-of-magnitude estimate, let’s round that volume off to the nearest
3
power of ten: 10 m . The shape of a tomato doesn’t follow linear dimensions, but since this is just an estimate, let’s pretend
3
that a tomato is an 0.1m by 0.1m by 0.1m cube, with a volume of 1 ⋅ 10 m . We can nd the total number of tomatoes by
−3 3
3 3
10 m
6
$2000
= $0.002 per tomato. That means that transportation really doesn’t contribute very much to the cost of a tomato.
10 tomatoes
Approximating the shape of a tomato as a cube is an example of another general strategy for making order-of-magnitude
estimates.
Key Points
Scientific notation means writing a number in terms of a product of something from 1 to 10 and something else that is a power
of 10.
In scientific notation all numbers are written in the form of a ⋅ 10 (a times ten raised to the power of b).
b
Each consecutive exponent number is ten times bigger than the previous one; negative exponents are used for small numbers.
When a sequence of calculations subject to rounding error is made, these errors can accumulate and lead to the
misrepresentation of calculated values.
Increasing the number of digits allowed in a representation reduces the magnitude of possible round-off errors, but may not
always be feasible, especially when doing manual calculations.
The degree to which numbers are rounded off is relative to the purpose of calculations and the actual value.
Orders of magnitude are generally used to make very approximate comparisons and reflect very large differences.
In the field of science, it is often sufficient for an estimate to be within an order of magnitude of the value in question.
1.3.3 https://phys.libretexts.org/@go/page/14433
When estimating area or volume, you are much better off estimating linear dimensions and computing volume from those linear
dimensions.
Key Terms
exponent: The power to which a number, symbol or expression is to be raised. For example, the 3 in x . 3
Scientific notation: A method of writing, or of displaying real numbers as a decimal number between 1 and 10 followed by an
integer power of 10
approximation: An imprecise solution or result that is adequate for a defined purpose.
Order of Magnitude: The class of scale or magnitude of any amount, where each class contains values of a fixed ratio (most
often 10) to the class preceding it. For example, something that is 2 orders of magnitude larger is 100 times larger; something
that is 3 orders of magnitude larger is 1000 times larger; and something that is 6 orders of magnitude larger is one million times
larger, because 10 = 100, 10 = 1000, and 10 = one million
2 3 6
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1.4: Solving Physics Problems
Dimensional Analysis
Any physical quantity can be expressed as a product of a combination of the basic physical dimensions.
learning objectives
Calculate the conversion from one kind of dimension to another
Dimensions
The dimension of a physical quantity indicates how it relates to one of the seven basic quantities. These fundamental quantities are:
[M] Mass
[L] Length
[T] Time
[A] Current
[K] Temperature
[mol] Amount of a Substance
[cd] Luminous Intensity
As you can see, the symbol is enclosed in a pair of square brackets. This is often used to represent the dimension of individual basic
[L]
quantity. An example of the use of basic dimensions is speed, which has a dimension of 1 in length and -1 in time; [T]
= [LT
−1
] .
Any physical quantity can be expressed as a product of a combination of the basic physical dimensions.
Dimensional Analysis
Dimensional analysis is the practice of checking relations between physical quantities by identifying their dimensions. The
dimension of any physical quantity is the combination of the basic physical dimensions that compose it. Dimensional analysis is
based on the fact that physical law must be independent of the units used to measure the physical variables. It can be used to check
the plausibility of derived equations, computations and hypotheses.
Derived Dimensions
The dimensions of derived quantities may include few or all dimensions in individual basic quantities. In order to understand the
technique to write dimensions of a derived quantity, we consider the case of force. Force is defined as:
F = m⋅a (1.4.1)
F = [M][a] (1.4.2)
The dimension of acceleration, represented as [a], is itself a derived quantity being the ratio of velocity and time. In turn, velocity is
also a derived quantity, being ratio of length and time.
−1
F = [M][a] = [M][vT ] (1.4.3)
−1 −1 −2
F = [M][LT T ] = [MLT ] (1.4.4)
Dimensional Conversion
In practice, one might need to convert from one kind of dimension to another. For common conversions, you might already know
how to convert off the top of your head. But for less common ones, it is helpful to know how to find the conversion factor:
Q = n1 u1 = n2 u2 (1.4.5)
where n represents the amount per u dimensions. You can then use ratios to figure out the conversion:
u2
n2 = ⋅ n1 (1.4.6)
u1
1.4.1 https://phys.libretexts.org/@go/page/14434
Trigonometry
Trigonometry is central to the use of free body diagrams, which help visually represent difficult physics problems.
learning objectives
Explain why trigonometry is useful in determining horizontal and vertical components of forces
Free Body Diagram: The rod is hinged from a wall and is held with the help of a string.
A rod ‘AB’ is hinged at ‘A’ from a wall and is held still with the help of a string, as shown in. This exercise involves drawing the
free body diagram. To make the problem easier, the force F will be expressed in terms of its horizontal and vertical components.
Removing all other elements from the image helps produce the finished free body diagram.
1.4.2 https://phys.libretexts.org/@go/page/14434
learning objectives
Construct a free-body diagram for a physical scenario
In physics, most problems are solved much more easily when a free body diagram is used. This uses geometry and vectors to
visually represent to problem, and trigonometry is also used in determining horizontal and vertical components of forces and
objects.
Purpose: Free body diagrams are very helpful in visually identifying which components are unknown, where the moments are
applied, and help analyze a problem, whether static or dynamic.
Do not include:
Do not show bodies other than the body of interest.
Do not show forces exerted by the body.
Internal forces acting on various parts of the body by other parts of the body.
Any velocity or acceleration is left out.
1.4.3 https://phys.libretexts.org/@go/page/14434
How To Solve Any Physics Problem
How To Solve Any Physics Problem: Learn five simple steps in five minutes! In this episode we cover the most effective
problem-solving method I’ve encountered and call upon some fuzzy friends to help us remember the steps.
Free Body Diagram: Use this figure to work through the example problem.
1.4.4 https://phys.libretexts.org/@go/page/14434
Key Points
Dimensional analysis is the practice of checking relations amount physical quantities by identifying their dimensions.
It is common to be faced with a problem that uses different dimensions to express the same basic quantity. The following
u2
equation can be used to find the conversion factor between the two derived dimensions: n = × n .
2
u1
1
Dimensional analysis can also be used as a simple check to computations, theories and hypotheses.
It is important to identify the problem and the unknowns and draw them in a free body diagram.
The laws of cosine and sine can be used to determine the vertical and horizontal components of the different elements of the
diagram.
Free body diagrams use geometry and vectors to visually represent physics problems.
A free body diagram lets you visually isolate the problem you are trying to solve, and simplify it into simple geometry and
trigonometry.
When drawing these diagrams, it is helpful to only draw the body it self, and the forces acting on it.
Drawing other objects and internal forces can condense the diagram and cause it to be less helpful.
Key Terms
dimension: A measure of spatial extent in a particular direction, such as height, width or breadth, or depth.
trigonometry: The branch of mathematics that deals with the relationships between the sides and the angles of triangles and the
calculations based on them, particularly the trigonometric functions.
static: Fixed in place; having no motion.
dynamic: Changing; active; in motion.
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CHAPTER OVERVIEW
10: Fluids
10.1: Introduction
10.2: Density and Pressure
10.3: Archimedes’ Principle
10.4: Cohesion and Adhesion
10.5: Fluids in Motion
10.6: Deformation of Solids
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1
10.1: Introduction
learning objectives
Assess the distinguishing characteristics of the four states of matter
There are a number of properties that can be observed in a material that identify what state the matter is in — solid, liquid, gas, or
plasma.
Solids
Solid: Solids are in a state of matter that maintains a fixed volume and shape.
A solid is in a state of matter that maintains a fixed volume and shape. A solid’s particles fit closely together. The forces between
the particles are so strong that the particles can not move freely; they can only vibrate. This causes a solid to be a stable, non-
compressible shape with definite volume.
Liquids
Liquid: Liquids maintain a fixed volume, but their shape will mold to the shape of the container they are being held in.
A liquid maintains a fixed volume, but its shape will mold to the shape of the container it is being held in. In, you can see that even
though the liquid’s shape is determined by the container, it has a free surface that is not controlled by the container. The particles
are close together but not as close as in solids; they are still able to move around, which causes the liquid to flow. Liquids usually
have a higher volume than their solid counterparts, except for certain molecules, such as H2O (water).
Gases
Gas: The particles are much farther from each other, usually a farther length than the size of the particles, and move a lot
shows a gas, whose particles move around a lot and are much farther apart from each other, usually farther apart than the diameter
of the particles themselves. The gas behaves like a liquid; the particles are moving but are still attracted to each other, so they still
flow. Unlike a solid or a liquid, the gas will try to fill whatever container it is in, adapting its volume accordingly.
Plasma
Plasma is a gas that has been ionized. That is to say, sufficient energy has been supplied to the gas such that the electrons have
enough energy to escape their atoms or molecules. Plasma contains ions and electrons that can move around freely. Matter in the
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plasma state has variable volume and shape. Plasma is the most common form of visible matter in the universe. Lightning, sparks,
neon lights, and the sun are all examples of matter in the plasma state.
Plasma: Matter in the plasma state has variable volume and shape, but as well as neutral atoms, it contains a significant number of
ions and electrons, both of which can move around freely.
Phase Transitions
shows the different states of matter and how they can change from one to another as a function of enthalpy and pressure and
temperature changes. A solid can change to a liquid with an enthalpy increase. The process of a liquid going to a solid is known as
melting. A liquid can change into a gas when it hits its boiling point or can even enter a plasma state if the enthalpy is increased
enough. When the enthalpy is lowered, a liquid can transform into a solid through freezing. Sometimes, a solid can skip freezing
and go directly to a gaseous state; this process is called sublimation.
States of Matter: This figure illustrates the relationship between the enthalpy of a system and the state of matter that the system is
in.
What is a Fluid?
A fluid is a substance that continually deforms (flows) under an applied shear stress.
learning objectives
Explain the properties of substances under an applied shear stress
A fluid is a substance that continually deforms (flows) under an applied shear stress. Fluids are a subset of the states of matter and
include three of the four states—liquids, gases, and plasma (shown in ).
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Four Fundamental States of Matter: Four fundamental states of matter: 1) top left corner corresponds to solid; 2) top right corner
corresponds to liquid; 3) bottom left corner corresponds to gas; 4) bottom right corner corresponds to plasma.
Liquids form a free surface (i.e., a surface not created by the container) while gases do not. The distinction between solids and fluid
is not entirely obvious. The distinction is made by evaluating the viscosity of the substance. Silly Putty can be considered to behave
like a solid or a fluid, depending on the time period over which it is observed. It is best described as a viscoelastic fluid.
Fluids display properties such as:
a) not resisting deformation or resisting it only lightly (viscosity), and
b) the ability to flow (also described as the ability to take on the shape of the container).
This also means that all fluids have the property of fluidity. These properties are typically a function of their inability to support a
shear stress in static equilibrium.
Solids can be subjected to shear stresses, and normal stresses—both compressive and tensile. In contrast, ideal fluids can only be
subjected to normal, compressive stress (called pressure). Real fluids display viscosity and so are capable of being subjected to low
levels of shear stress.
Although the term fluid includes both the liquid and gas phases, it is also commonly used as a synonym for liquid, with no
implication that gas could also be present. For example, “brake fluid” is hydraulic oil and will not perform its required function if
there is gas in it. This colloquial usage of the term is also common in the fields of medicine and nutrition (e.g., “take plenty of
fluids”).
Key Points
Solids are non-compressible and have constant volume and constant shape.
Liquids are non-compressible and have constant volume but can change shape. A liquid’s shape is dictated by the shape of the
container it is in.
Gases do not have a constant volume or shape; they not only take the shape of the container they are in, they try to fill the entire
container.
Matter in the plasma state has variable volume and shape. Plasma contains ions and electrons, both of which can move around
freely.
Fluids are a subset of the states of matter and include liquids, gases, and plasma.
Fluids display properties such as: not resisting deformation or resisting it only lightly ( viscosity ), and the ability to flow (also
described as the ability to take on the shape of the container).
Ideal fluids can only be subjected to normal, compressive stress which is called pressure. Real fluids display viscosity and thus
are capable of being subjected to low levels of shear stress.
Key Items
plasma: a state of matter consisting of partially ionized gas
enthalpy: the total amount of energy in a system, including both the internal energy and the energy needed to displace its
environment
sublimation: the transition of a substance from the solid phase directly to the vapor state such that it does not pass through the
intermediate, liquid phase
fluidity: A measure of the extent to which something is fluid. The reciprocal of its viscosity.
viscosity: A quantity expressing the magnitude of internal friction in a fluid, as measured by the force per unit area resisting
uniform flow.
shear stress: The component of stress that causes parallel layers of a material to move relative to each other in their own
planes.
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10.2: Density and Pressure
learning objectives
Identify factors that determine the pressure exerted by the gas
Pressure is an important physical quantity—it plays an essential role in topics ranging from thermodynamics to solid and fluid
mechanics. As a scalar physical quantity (having magnitude but no direction), pressure is defined as the force per unit area applied
perpendicular to the surface to which it is applied. Pressure can be expressed in a number of units depending on the context of use.
Pressure and Pascal’s Principle: A brief introduction to pressure and Pascal’s Principle, including hydraulics.
where p is pressure, F is the force acting perpendicular to the surface to which this force is applied, and A is the area of the
surface. Any object that possesses weight, whether at rest or not, exerts a pressure upon the surface with which it is in contact. The
magnitude of the pressure exerted by an object on a given surface is equal to its weight acting in the direction perpendicular to that
surface, divided by the total surface area of contact between the object and the surface. shows the graphical representations and
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corresponding mathematical expressions for the case in which a force acts perpendicular to the surface of contact, as well as the
case in which a force acts at angle θ relative to the surface.
Representation of Pressure: This image shows the graphical representations and corresponding mathematical expressions for the
case in which a force acts perpendicular to the surface of contact, as well as the case in which a force acts at angle θ relative to the
surface.
Pressure as a Function of Surface Area: Pressure can be increased by either increasing the force or by decreasing the area or can
oppositely be decreased by either decreasing the force or increasing the area.
A good illustration of this is the reason a sharp knife is far more effective for cutting than a blunt knife. The same force applied by
a sharp knife with a smaller area of contact will exert a much greater pressure than a blunt knife having a considerably larger area
of contact. Similarly, a person standing on one leg on a trampoline causes a greater displacement of the trampoline than that same
person standing on the same trampoline using two legs—not because the individual exerts a larger force when standing on one leg,
but because the area upon which this force is exerted is decreased, thus increasing the pressure on the trampoline. Alternatively, an
object having a weight larger than another object of the same dimensionality and area of contact with a given surface will exert a
greater pressure on that surface due to an increase in force. Finally, when considering a given force of constant magnitude acting on
a constant area of a given surface, the pressure exerted by that force on that surface will be greater the larger the angle of that force
as it acts upon the surface, reaching a maximum when that force acts perpendicular to the surface.
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where n is the number of gas molecules, R is the ideal gas constant (R = 8.314 J mol-1 K-1), T is the temperature of the gas, and V
is the volume of the container.
The pressure exerted by the gas can be increased by: increasing the number of collisions of gas molecules per unit time by
increasing the number of gas molecules; increasing the kinetic energy of the gas by increasing the temperature; or decreasing the
volume of the container. offers a representation of the ideal gas law, as well as the effect of varying the equation parameters on the
gas pressure. Another common type of pressure is that exerted by a static liquid or hydrostatic pressure. Hydrostatic pressure is
most easily addressed by treating the liquid as a continuous distribution of matter, and may be considered a measure of energy per
unit volume or energy density. We will further discuss hydrostatic pressure in other sections.
Pressure of an Ideal Gas: This image is a representation of the ideal gas law, as well as the effect of varying the equation
parameters on the gas pressure.
learning obectives
Identify factors that determine the pressure exerted by static liquids and gases
Pressure is defined in simplest terms as force per unit area. However, when dealing with pressures exerted by gases and liquids, it is
most convenient to approach pressure as a measure of energy per unit volume by means of the definition of work (W = F·d). The
derivation of pressure as a measure of energy per unit volume from its definition as force per unit area is given in. Since, for gases
and liquids, the force acting on a system contributing to pressure does not act on a specific point or particular surface, but rather as
a distribution of force, analyzing pressure as a measure of energy per unit volume is more appropriate. For liquids and gases at rest,
the pressure of the liquid or gas at any point within the medium is called the hydrostatic pressure. At any such point within a
medium, the pressure is the same in all directions, as if the pressure was not the same in all directions, the fluid, whether it is a gas
or liquid, would not be static. Note that the following discussion and expressions pertain only to incompressible fluids at static
equilibrium.
Energy per Unit Volume: This equation is the derivation of pressure as a measure of energy per unit volume from its definition as
force per unit area.
The pressure exerted by a static liquid depends only on the depth, density of the liquid, and the acceleration due to gravity. gives
the expression for pressure as a function of depth within an incompressible, static liquid as well as the derivation of this equation
from the definition of pressure as a measure of energy per unit volume (ρ is the density of the gas, g is the acceleration due to
gravity, and h is the depth within the liquid). For any given liquid with constant density throughout, pressure increases with
increasing depth. For example, a person under water at a depth of h1 will experience half the pressure as a person under water at a
depth of h = 2h . For many liquids, the density can be assumed to be nearly constant throughout the volume of the liquid and, for
2 1
virtually all practical applications, so can the acceleration due to gravity (g = 9.81 m/s2). As a result, pressure within a liquid is
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therefore a function of depth only, with the pressure increasing at a linear rate with respect to increasing depth. In practical
applications involving calculation of pressure as a function of depth, an important distinction must be made as to whether the
absolute or relative pressure within a liquid is desired. Equation 2 by itself gives the pressure exerted by a liquid relative to
atmospheric pressure, yet if the absolute pressure is desired, the atmospheric pressure must then be added to the pressure exerted by
the liquid alone.
Pressure as Energy per Unit Volume: This equation gives the expression for pressure as a function of depth within an
incompressible, static liquid as well as the derivation of this equation from the definition of pressure as a measure of energy per
unit volume (ρ is the density of the gas, g is the acceleration due to gravity, and h is the depth within the liquid).
When analyzing pressure within gases, a slightly different approach must be taken as, by the nature of gases, the force contributing
to pressure arises from the average number of gas molecules occupying a certain point within the gas per unit time. Thus the force
contributing to the pressure of a gas within the medium is not a continuous distribution as for liquids and the barometric equation
given in must be utilized to determine the pressure exerted by the gas at a certain depth (or height) within the gas (p0 is the pressure
at h = 0, M is the mass of a single molecule of gas, g is the acceleration due to gravity, k is the Boltzmann constant, T is the
temperature of the gas, and h is the height or depth within the gas). Equation 3 assumes that the gas is incompressible and that the
pressure is hydrostatic.
Pressure within a gas: The force contributing to the pressure of a gas within the medium is not a continuous distribution as for
liquids and the barometric equation given in this figure must be utilized to determine the pressure exerted by the gas at a certain
depth (or height) within the gas (p0 is the pressure at h = 0, M is the mass of a single molecule of gas, g is the acceleration due to
gravity, k is the Boltzmann constant, T is the temperature of the gas, and h is the height or depth within the gas)
Static Equilibrium
Any region or point, or any static object within a static fluid is in static equilibrium where all forces and torques are equal to zero.
learning obectives
Identify required conditions for a fluid to be in rest
Static equilibrium is a particular state of a physical system. It is qualitatively described by an object at rest and by the sum of all
forces, with the sum of all torques acting on that object being equal to zero. Static objects are in static equilibrium, with the net
force and net torque acting on that object being equal to zero; otherwise there would be a driving mechanism for that object to
undergo movement in space. The analysis and study of objects in static equilibrium and the forces and torques acting on them is
called statics—a subtopic of mechanics. Statics is particularly important in the design of static and load bearing structures. As it
pertains to fluidics, static equilibrium concerns the forces acting on a static object within a fluid medium.
Fluids
For a fluid at rest, the conditions for static equilibrium must be met at any point within the fluid medium. Therefore, the sum of the
forces and torques at any point within the static liquid or gas must be zero. Similarly, the sum of the forces and torques of an object
at rest within a static fluid medium must also be zero. In considering a stationary object within a liquid medium at rest, the forces
acting at any point in time and at any point in space within the medium must be analyzed. For a stationary object within a static
liquid, there are no torques acting on the object so the sum of the torques for such a system is immediately zero; it need not concern
analysis since the torque condition for equilibrium is fulfilled.
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Density
At any point in space within a static fluid, the sum of the acting forces must be zero; otherwise the condition for static equilibrium
would not be met. In analyzing such a simple system, consider a rectangular region within the fluid medium with density ρL (same
density as the fluid medium), width w, length l, and height h, as shown in. Next, the forces acting on this region within the medium
are taken into account. First, the region has a force of gravity acting downwards (its weight) equal to its density object, times its
volume of the object, times the acceleration due to gravity. The downward force acting on this region due to the fluid above the
region is equal to the pressure times the area of contact. Similarly, there is an upward force acting on this region due to the fluid
below the region equal to the pressure times the area of contact. For static equilibrium to be achieved, the sum of these forces must
be zero, as shown in. Thus for any region within a fluid, in order to achieve static equilibrium, the pressure from the fluid below the
region must be greater than the pressure from the fluid above by the weight of the region. This force which counteracts the weight
of a region or object within a static fluid is called the buoyant force (or buoyancy).
Static Equilibrium of a Region Within a Fluid: This figure shows the equations for static equilibrium of a region within a fluid.
Region Within a Static Fluid: This figure is a free body diagram of a region within a static fluid.
In the case on an object at stationary equilibrium within a static fluid, the sum of the forces acting on that object must be zero. As
previously discussed, there are two downward acting forces, one being the weight of the object and the other being the force
exerted by the pressure from the fluid above the object. At the same time, there is an upwards force exerted by the pressure from
the fluid below the object, which includes the buoyant force. shows how the calculation of the forces acting on a stationary object
within a static fluid would change from those presented in if an object having a density ρS different from that of the fluid medium is
surrounded by the fluid. The appearance of a buoyant force in static fluids is due to the fact that pressure within the fluid changes
as depth changes. The analysis presented above can furthermore be extended to much more complicated systems involving
complex objects and diverse materials.
Pascal’s Principle
Pascal’s Principle states that pressure is transmitted and undiminished in a closed static fluid.
learning obectives
Apply Pascal’s Principle to describe pressure behavior in static fluids
Pascal’s Principle
Pascal’s Principle (or Pascal’s Law ) applies to static fluids and takes advantage of the height dependency of pressure in static
fluids. Named after French mathematician Blaise Pascal, who established this important relationship, Pascal’s Principle can be used
to exploit pressure of a static liquid as a measure of energy per unit volume to perform work in applications such as hydraulic
presses. Qualitatively, Pascal’s Principle states that pressure is transmitted undiminished in an enclosed static liquid. Quantitatively,
Pascal’s Law is derived from the expression for determining the pressure at a given height (or depth) within a fluid and is defined
by Pascal’s Principle:
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AP Physics 2 - Pressure and Pascal's Principle
Pressure and Pascal’s Principle: A brief introduction to pressure and Pascal’s Principle, including hydraulics.
where p1 is the external applied pressure, ρ is the density of the fluid, Δh is the difference in height of the static liquid, and g is the
acceleration due to gravity. Pascal’s Law explicitly determines the pressure difference between two different heights (or depths)
within a static liquid. As, by Pascal’s Law, a change in pressure is linearly proportional to a change in height within an
incompressible, static liquid of constant density, doubling the height between the two points of reference will double the change of
pressure, while halving the height between the two points will half the change in pressure.
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weight of the static fluid in the bottle. This example shows how, through Pascal’s Principle, the force exerted by a static fluid in a
closed system can be multiplied by changing the height and the surface area of contact.
Pressure Applied to a Hydrostatic Fluid: A downwards force of 10 N is applied to a bottle filled with a static liquid of constant
density ρ at the spout of cross-sectional area of 5 cm2, yielding an applied pressure of 2 N/cm2.
Hydraulic Press Diagrams: Two different types of hydraulic press configurations, the first in which there is no difference in
height of the static liquid and the second in which there is a difference in height Δh of the static liquid.
learning obectives
Explain the relationship among absolute pressure, gauge pressure, and atmospheric pressure
Atmospheric Pressure
An important distinction must be made as to the type of pressure quantity being used when dealing with pressure measurements
and calculations. Atmospheric pressure is the magnitude of pressure in a system due to the atmosphere, such as the pressure exerted
by air molecules (a static fluid ) on the surface of the earth at a given elevation. In most measurements and calculations, the
atmospheric pressure is considered to be constant at 1 atm or 101,325 Pa, which is the atmospheric pressure under standard
conditions at sea level.
Atmospheric pressure is due to the force of the molecules in the atmosphere and is a case of hydrostatic pressure. Depending on the
altitude relative to sea level, the actual atmospheric pressure will be less at higher altitudes and more at lower altitudes as the
weight of air molecules in the immediate atmosphere changes, thus changing the effective atmospheric pressure. Atmospheric
pressure is a measure of absolute pressure and can be affected by the temperature and air composition of the atmosphere but can
generally be accurately approximated to be around standard atmospheric pressure of 101,325 Pa. Within the majority of earth’s
atmosphere, pressure varies with height according to. In this equation p0 is the pressure at sea level (101,325 Pa), g is the
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acceleration due to gravity, M is the mass of a single molecule of air, R is the universal gas constant, T0 is the standard temperature
at sea level, and h is the height relative to sea level.
Gauge Pressure
For most applications, particularly those involving pressure measurements, it is more practical to use gauge pressure than absolute
pressure as a unit of measurement. Gauge pressure is a relative pressure measurement which measures pressure relative to
atmospheric pressure and is defined as the absolute pressure minus the atmospheric pressure. Most pressure measuring equipment
give the pressure of a system in terms of gauge pressure as opposed to absolute pressure. For example, tire pressure and blood
pressure are gauge pressures by convention, while atmospheric pressures, deep vacuum pressures, and altimeter pressures must be
absolute.
For most working fluids where a fluid exists in a closed system, gauge pressure measurement prevails. Pressure instruments
connected to the system will indicate pressures relative to the current atmospheric pressure. The situation changes when extreme
vacuum pressures are measured; absolute pressures are typically used instead.
To find the absolute pressure of a system, the atmospheric pressure must then be added to the gauge pressure. While gauge pressure
is very useful in practical pressure measurements, most calculations involving pressure, such as the ideal gas law, require pressure
values in terms of absolute pressures and thus require gauge pressures to be converted to absolute pressures.
learning obectives
Compare design and operation of aneroid and hydrostatic based barometers
Gauge Pressure
In practice, pressure is most often measured in terms of gauge pressure. Gauge pressure is the pressure of a system above
atmospheric pressure. Since atmospheric pressure is mostly constant with little variation near sea level, where most practical
pressure measurements are taken, it is assumed to be approximately 101,325 Pa. Modern pressure measuring devices sometimes
have incorporated mechanisms to account for changes in atmospheric pressure due to elevation changes. Gauge pressure is much
more convenient than absolute pressure for practical measurements and is widely used as an established measure of pressure.
However, it is important to determine whether it is necessary to use absolute (gauge plus atmospheric) pressure for calculations, as
is often the case for most calculations, such as those involving the ideal gas law. Pressure measurements have been accurately taken
since the mid-1600s with the invention of the traditional barometer. Barometers are devices used to measure pressure and were
initially used to measure atmospheric pressure.
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mercury, mercury is often a better choice for fabricating accurate hydrostatic barometers. The density of mercury is much higher
than that of water, thus allowing for higher accuracy of measurements and the ability to fabricate more compact hydrostatic
barometers. In theory, a hydrostatic barometer can be placed in a closed system to measure the absolute pressure and the gauge
pressure of the system by subtracting the atmospheric pressure.
Aneroid Barometer
Another type of barometer is the aneroid barometer, which consists of a small, flexible sealed metal box called an aneroid cell. The
aneroid cell is made from beryllium-copper alloy and is partially evacuated. A stiff spring prevents the aneroid cell from
collapsing. Small changes in external air pressure cause the cell to expand or contract. This expansion and contraction is amplified
by mechanical mechanisms to give a pressure reading. Such pressure measuring devices are more practical than hydrostatic
barometers for measuring system pressures. Many modern pressure measuring devices are pre-engineered to output gauge pressure
measurements. While the aneroid barometer is the underlying mechanism behind many modern pressure measuring devices,
pressure can also be measured using more advanced measuring mechanisms.
Hydrostatic Column Barometer: The concept of determining pressure using the fluid height in a hydrostatic column barometer
Variation of Pressure with Height: The density of the liquid is p, g is the acceleration due to gravity, and h is the height of the
fluid in the barometer column.
learning obectives
Explain role played by pressure in the circulatory and respiratory systems
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Approximation for Mean Arterial Pressure: In practice, the mean arterial pressure (MAP) can be approximated from easily
obtainable blood pressure measurements.
The mean arterial pressure (MAP) is the average pressure over a cardiac cycle and is determined by, where CO is the cardiac
outputs, SVR is the systemic vascular resistance, and CVP is the central venous pressure (CVP). In practice, the mean arterial
pressure (MAP) can be approximated from easily obtainable blood pressure measurements in, where Psys is the measured systolic
pressure and Pdias is the measured diastolic pressure. One particularly common and dangerous circulatory system condition is
partial blockage of blood vessels due to a number of factors, such as plaque build-up from high cholesterol, which results in a
reduction of the effective blood vessel cross-sectional diameter and a corresponding reduction in blood flow rate and thus an
increase in blood pressure to restore normal blood flow according to Poiseuille’s Law.
Equation for Mean Arterial Pressure: The mean arterial pressure (MAP) is the average pressure over a cardiac cycle and is
determined this equation, where CO is the cardiac outputs, SVR is the systemic vascular resistance, and CVP is the central venous
pressure (CVP).
Key Points
Pressure is a scalar quantity defined as force per unit area. Pressure only concerns the force component perpendicular to the
surface upon which it acts, thus if the force acts at an angle, the force component along the direction perpendicular to the
surface must be used to calculate pressure.
The pressure exerted on a surface by an object increases as the weight of the object increases or the surface area of contact
decreases. Alternatively the pressure exerted decreases as the weight of the object decreases or the surface area of contact
increases.
Pressure exerted by ideal gases in confined containers is due to the average number of collisions of gas molecules with the
container walls per unit time. As such, pressure depends on the amount of gas (in number of molecules), its temperature, and
the volume of the container.
Hydrostatic pressure refers to the pressure exerted by a fluid (gas or liquid) at any point in space within that fluid, assuming that
the fluid is incompressible and at rest.
Pressure within a liquid depends only on the density of the liquid, the acceleration due to gravity, and the depth within the
liquid. The pressure exerted by such a static liquid increases linearly with increasing depth.
Pressure within a gas depends on the temperature of the gas, the mass of a single molecule of the gas, the acceleration due to
gravity, and the height (or depth) within the gas.
Hydrostatic balance is the term used for a region or stationary object within a static fluid which is at static equilibrium, and for
which the sum of all forces and sum of all torques is equal to zero.
A region or static object within a stationary fluid experiences downward forces due to the weight of the region or object, and the
pressure exerted from the fluid above the region or object, as well as an upward force due to the pressure exerted from the fluid
10.2.10 https://phys.libretexts.org/@go/page/14495
below the region or object.
For a region or static object within a static fluid, the downward force due to the weight of the region or object is counteracted by
the upward buoyant force, which is equal to the weight of the fluid displaced by the region or object.
Pascal’s Principle is used to quantitatively relate the pressure at two points in an incompressible, static fluid. It states that
pressure is transmitted, undiminished, in a closed static fluid.
The total pressure at any point within an incompressible, static fluid is equal to the sum of the applied pressure at any point in
that fluid and the hydrostatic pressure change due to a difference in height within that fluid.
Through the application of Pascal’s Principle, a static liquid can be utilized to generate a large output force using a much
smaller input force, yielding important devices such as hydraulic presses.
Atmospheric pressure is a measure of absolute pressure and is due to the weight of the air molecules above a certain height
relative to sea level, increasing with decreasing altitude and decreasing with increasing altitude.
Gauge pressure is the additional pressure in a system relative to atmospheric pressure. It is a convenient pressure measurement
for most practical applications.
While gauge pressure is more convenient for practical measurements, absolute pressure is necessary for most pressure
calculations, thus the atmospheric pressure must be added to the gauge pressure for calculations.
Gauge pressure is the pressure of a system above atmospheric pressure, which must be converted to absolute pressure for most
calculations.
The barometer is a device which uses hydrostatic fluids to directly determine atmospheric pressure and may be used to
indirectly measure the gauge pressure of systems.
The hydrostatic column barometer uses a liquid like water or mercury for functionality, while the aneroid barometer uses an
evacuated flexible metal cell.
Pressure, along with the potential for work arising from differences in pressure, plays an essential role in the functionality of
several critical bodily functions and systems necessary for survival.
The circulatory system relies on pressure differences for circulating blood, along with oxygen, necessary nutrients, and waste
products throughout the body.
Respiration is made possible as a result of pressure differences between the thoracic cavity, the lungs, and the environment and
is largely regulated by movement of the diaphragm.
Key Terms
ideal gas: Theoretical gas characterized by random motion whose individual molecules do not interact with one another and are
chemically inert.
kinetic energy: The energy associated with a moving particle or object having a certain mass.
incompressible: Unable to be compressed or condensed.
static equilibrium: the physical state in which all components of a system are at rest and the net force is equal to zero
throughout the system
Buoyancy: The power of supporting a body so that it floats; upward pressure exerted by the fluid in which a body is immersed.
torque: Something that produces or tends to produce torsion or rotation; the moment of a force or system of forces tending to
cause rotation.
equilibrium: A state of rest or balance due to the equal action of opposing forces.
hydraulic press: Device that uses a hydraulic cylinder (closed static fluid) to generate a compressive force.
Gauge Pressure: The pressure of a system above atmospheric pressure.
Torr: A unit of pressure equal to one millimeter of mercury (760 torr = 101,325 Pa).
Aneroid Barometer: A device for measuring pressure, often specially calibrated for use as an altimeter, consisting of a box or
chamber partially exhausted of air, having an elastic top and a pointer to indicate the degree of compression of the top caused by
the external air.
Thoracic Cavity: A hollow place or space, or a potential space, within the body or one of its organs.
Poiseuille’s Law: The law that the velocity of a liquid flowing through a capillary is directly proportional to the pressure of the
liquid and the fourth power of the radius of the capillary and is inversely proportional to the viscosity of the liquid and the
length of the capillary.
Alveoli: Small air sacs or cavities in the lung that give the tissue a honeycomb appearance and expand its surface area for the
exchange of oxygen and carbon dioxide.
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10.3: Archimedes’ Principle
learning objectives
Calculate the direction of the buoyancy force
When you rise from soaking in a warm bath, your arms may feel strangely heavy. This effect is due to the loss of the buoyant
support of the water. What creates this buoyant force ? Why is it that some things float and others do not? Do objects that sink get
any support at all from the fluid? Is your body buoyed by the atmosphere, or are only helium balloons affected?
P1 = h1 ρg, (10.3.1)
s2
is the gravitational acceleration. The magnitude of the force on the top surface
is:
F1 = P1 A = h1 ρgA. (10.3.2)
This force points downwards. Similarly, the force on the bottom surface is:
F2 = P2 A = h2 ρgA (10.3.3)
and points upwards. Because it is cylindrical, the net force on the object’s sides is zero—the forces on different parts of the surface
oppose each other and cancel exactly. Thus, the net upward force on the cylinder due to the fluid is:
FB = F2 − F1 = ρgA(h2 − h1 ) (10.3.4)
Archimedes principle: The buoyant force on the ship (a) is equal to the weight of the water displaced by the ship—shown as the
dashed region in (b).
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The principle can be stated as a formula:
FB = wf l (10.3.5)
The reasoning behind the Archimedes principle is that the buoyancy force on an object depends on the pressure exerted by the fluid
on its submerged surface. Imagine that we replace the submerged part of the object with the fluid in which it is contained, as in (b).
The buoyancy force on this amount of fluid must be the same as on the original object (the ship). However, we also know that the
buoyancy force on the fluid must be equal to its weight, as the fluid does not sink in itself. Therefore, the buoyancy force on the
original object is equal to the weight of the “displaced fluid” (in this case, the water inside the dashed region (b)).
The Archimedes principle is valid for any fluid—not only liquids (such as water) but also gases (such as air). We will explore this
further as we discuss applications of the principle in subsequent sections.
Archimedes’ Principle – Simple Example: We use Archimedes’ Principle to determine the number of penguins an ice float can
dryly support.
Complete Submersion
The buoyancy force on a completely submerged object of volume is F B = Vρg .
learning objectivies
Identify factors determining the buoyancy force on a completely submerged object
The Archimedes principle is easiest to understand and apply in the case of entirely submersed objects. In this section we discuss a
few relevant examples. In general, the buoyancy force on a completely submerged object is given by the formula:
10.3.2 https://phys.libretexts.org/@go/page/14496
FB = Vρg, (10.3.6)
where V is the volume of the object, ρ is the density of the fluid, and g is gravitational acceleration. This follows immediately from
the Archimedes’ principle, and the fact that the object is completely submerged (and so the volume of the fluid displaced is just the
volume of the object).
Cylinder
In the previous section, we calculated the buoyancy force on a cylinder (shown in ) by considering the force exerted on each of the
cylinder’s sides. Now, we’ll calculate this force using Archimedes’ principle. The buoyancy force on the cylinder is equal to the
weight of the displaced fluid. This weight is equal to the mass of the displaced fluid multiplied by the gravitational acceleration:
Buoyant force: The fluid pushes on all sides of a submerged object. However, because pressure increases with depth, the upward
push on the bottom surface (F2) is greater than the downward push on the top surface (F1). Therefore, the net buoyant force is
always upwards.
FB = wf l = mf l g (10.3.7)
The mass of the displaced fluid is equal to its volume multiplied by its density:
mf l = Vf l ρ (10.3.8)
However (and this is the crucial point), the cylinder is entirely submerged, so the volume of the displaced fluid is just the volume
of the cylinder (see ), and:
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Archimedes principle: The volume of the fluid displaced (b) is the same as the volume of the original cylinder (a).
mf l = Vf l ρ = Vcylinder ρ. (10.3.9)
The volume of a cylinder is the area of its base multiplied by its height, or in our case:
Vcylinder = A(h2 − h1 ). (10.3.10)
This is the same result obtained in the previous section by considering the force due to the pressure exerted by the fluid.
Helium Airship
Consider the USS Macon, a helium-filled airship (shown in ). Its envelope (the “balloon”) contained 184,059.5 cubic meters of
helium. Ignoring the small volume of the gondola, what was the buoyancy force on this airship? If the airship weighed 108,000 kg,
how much cargo could it carry? Assume the density of air is 1.225 kg per meter cubed. The buoyancy force on an airship is due to
the air in which it is immersed. Although we don’t know the exact shape of the airship, we know its volume and the density of the
air, and thus we can calculate the buoyancy force:
To find the cargo capacity of the airship, we subtract the weight of the airship from the buoyancy force:
m
6 5 6
Fcargo = FB − mg = 2.21 × 10 N − 1.08 × 10 kg × 9.81 = 1.15 × 10 N (10.3.13)
2
s
Flotation
An object floats if the buoyancy force exerted on it by the fluid balances its weight.
learning objectives
Express the relationship between the buoyancy force and the weight for a floating object
Why do some objects float, but others don’t? If you put a metal coin into a glass of water it will sink. But most ships are built of
metal, and they float. So how is this possible?
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Condition for Flotation
An object will float if the buoyancy force exerted on it by the fluid balances its weight, i.e. if FB=mgFB=mg.
But the Archimedes principle states that the buoyant force is the weight of the fluid displaced. So, for a floating object on a liquid,
the weight of the displaced liquid is the weight of the object. Thus, only in the special case of floating does the buoyant force acting
on an object equal the object’s weight. Consider a one-ton block of solid iron. As iron is nearly eight times denser than water, it
displaces only 1/8 ton of water when submerged, which is not enough to keep it afloat. Suppose the same iron block is reshaped
into a bowl. It still weighs one ton, but when it is put in water, it displaces a greater volume of water than when it was a block. The
deeper the iron bowl is immersed, the more water it displaces, and the greater the buoyant force acting on it. When the buoyant
force equals one ton, it will sink no further.
When any boat displaces a weight of water equal to its own weight, it floats. This is often called the “principle of flotation” where a
floating object displaces a weight of fluid equal to its own weight. Every ship, submarine, and dirigible must be designed to
displace a weight of fluid equal to its own weight. A 10,000-ton ship must be built wide enough to displace 10,000 tons of water
before it sinks too deep in the water. The same is true for vessels in air (as air is a fluid): A dirigible that weighs 100 tons displaces
at least 100 tons of air; if it displaces more, it rises; if it displaces less, it falls. If the dirigible displaces exactly its weight, it hovers
at a constant altitude.
Density and Submersion: An unloaded ship (a) floats higher in the water than a loaded ship (b).
Vsub Vf l
f raction submerged = = (10.3.15)
Vobj Vobj
The volume submerged equals the volume of fluid displaced, which we call Vf l . Now we can obtain the relationship between the
densities by substituting ρ = mV into the expression. This gives
mf l / ρf l
f raction submerged = (10.3.16)
mobj / ρ̄
obj
where ρ̄ is the average density of the object and ρflρflis the density of the fluid. Since the object floats, its mass and that of the
obj
displaced fluid are equal, and so they cancel from the equation, leaving
ρ̄ obj
f raction submerged = (10.3.17)
ρf l .
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between that of air and steel. To be more precise, the average density is defined as the total mass of an object divided by its total
volume: ρ̄ = .m
2. This formula makes sense only if the density of the object is smaller than the density of the fluid. Otherwise, the fraction
submerged becomes greater than one—a sign that the object does not float at all, but it sinks!
Key Points
The buoyancy force is caused by the pressure exerted by the fluid in which an object is immersed.
The buoyancy force always points upwards because the pressure of a fluid increases with depth.
You can calculate the buoyancy force either directly by computing the force exerted on each of the object’s surfaces, or
indirectly by finding the weight of the displaced fluid.
If an object is completely submerged, the volume of the fluid displaced is equal to the volume of the object.
The buoyancy force on hot-air balloons, dirigibles and other objects can be calculated by assuming that they are entirely
submerged in air.
The buoyancy force does not depend on the shape of the object, only on its volume.
The buoyancy force experienced by an object depends on its shape.
ρ̄
The fraction of an object’s volume that’s submerged is given by the ratio of its average density to that of the fluid: .
obj
ρ
fl
An object floats if the buoyancy force exerted on it by the fluid balances its weight.
Key Terms
buoyant force: An upward force exerted by a fluid that opposes the weight of an immersed object.
Archimedes principle: The buoyant force exerted on a body immersed in a fluid is equal to the weight of the fluid the body
displaces.
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10.4: Cohesion and Adhesion
learning objectives
Explain the phenomena of surface tension and capillary action
Attractive forces between molecules of the same type are called cohesive forces. Liquids can, for example, be kept in open
containers because cohesive forces hold the molecules together. Attractive forces between molecules of different types are called
adhesive forces. Such forces cause liquid drops to cling to window panes, for example. In this section we examine effects of
cohesive and adhesive forces in liquids.
Surface Tension
Surface tension is a contractive tendency of the surface of a liquid that allows it to resist an external force. It is shown, for example,
in the floating of some objects on the surface of water, even though they are denser than water, and in the ability of some insects
(e.g., water striders) to run on water’s surface. This property is caused by cohesion of similar molecules and is responsible for
many of the behaviors of liquids.
The cohesive forces among liquid molecules are responsible for the phenomenon of surface tension, as shown in. In the bulk of the
liquid, each molecule is pulled equally in every direction by neighboring liquid molecules, resulting in a net force of zero. The
molecules at the surface do not have other molecules on all sides of them and therefore are pulled inwards. This creates some
internal pressure and forces liquid surfaces to contract to the minimal area.
Capillary Action
Capillary action, or capillarity, is the ability of a liquid to flow in narrow spaces without the assistance of, and in opposition to,
external forces like gravity. The effect can be seen in the drawing-up of liquids between the hairs of a paint-brush, in a thin tube, in
porous materials such as paper, in some non-porous materials such as liquified carbon fiber, and in a cell. It occurs because of
intermolecular attractive forces between the liquid and solid surrounding surfaces. If the diameter of the tube is sufficiently small,
then the combination of surface tension (which is caused by cohesion within the liquid) and adhesive forces between the liquid and
the container act to lift the liquid.
With some pairs of materials, such as mercury and glass (see ), the intermolecular forces within the liquid exceed those between the
solid and the liquid, so a convex meniscus forms, and capillary action works in reverse.
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Capillarity: Capillary action of water compared to mercury, in each case with respect to glass
Key Points
Attractive forces between molecules of the same type are called cohesive forces.
Attractive forces between molecules of different types are called adhesive forces.
Surface tension is a contractive tendency of the surface of a liquid that allows it to resist an external force.
Capillary action, or capillarity, is the ability of a liquid to flow in narrow spaces without the assistance of, and in opposition to,
external forces such as gravity.
Key Term
Pressure: the amount of force that is applied over a given area divided by the size of that area
intermolecular: from one molecule to another; between molecules
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10.5: Fluids in Motion
learning objectives
Determine the flow rate based on velocity and area or elapsed time and justify the use of continuity in expressing properties
of a fluid and its motion
The flow rate of a fluid is the volume of fluid which passes through a surface in a given unit of time. It is usually represented by the
symbol Q.
Continuity Equation for Fluids: A brief introduction to the Continuity Equation for Fluids.
Flow Rate
Volumetric flow rate is defined as
Q = v × a, (10.5.1)
where Q is the flow rate, v is the velocity of the fluid, and a is the area of the cross section of the space the fluid is moving through.
Volumetric flow rate can also be found with
V
Q= (10.5.2)
t
where Q is the flow rate, V is the Volume of fluid, and t is elapsed time.
10.5.1 https://phys.libretexts.org/@go/page/14498
Continuity
The equation of continuity works under the assumption that the flow in will equal the flow out. This can be useful to solve for
many properties of the fluid and its motion:
Flow in = Flow out: Using the known properties of a fluid in one condition, we can use the continuity equation to solve for the
properties of the same fluid under other conditions.
Q1 = Q2 (10.5.3)
This can be expressed in many ways, for example: A v = A v . The equation of continuity applies to any incompressible fluid.
1 1 2 2
Since the fluid cannot be compressed, the amount of fluid which flows into a surface must equal the amount flowing out of the
surface.
Key Points
Flow rate can be expressed in either terms of cross sectional area and velocity, or volume and time.
Because liquids are incompressible, the rate of flow into an area must equal the rate of flow out of an area. This is known as the
equation of continuity.
The equation of continuity can show how much the speed of a liquid increases if it is forced to flow through a smaller area. For
example, if the area of a pipe is halved, the velocity of the fluid will double.
Although gases often behave as fluids, they are not incompressible the way liquids are and so the continuity equation does not
apply.
Key Terms
incompressible: Unable to be compressed or condensed.
continuity: Lack of interruption or disconnection; the quality of being continuous in space or time.
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10.6: Deformation of Solids
Learning objectives
Explain how length of an object is determined
Length
In geometric measurements, length is the longest dimension of an object. In other contexts “length” is the measured dimension of
an object. For example: it is possible to cut a length of a wire which is shorter than wire thickness. Length may be distinguished
from height, which is vertical extent, and width or breadth, which are the distance from side to side, measuring across the object at
right angles to the length.
Length is a measure of one dimension, whereas area is a measure of two dimensions (length squared) and volume is a measure of
three dimensions (length cubed). In most systems of measurement, the unit of length is a fundamental unit, from which other units
are defined.
After Albert Einstein’s Special Relativity Theory, length can no longer be thought of being constant in all reference frames. Thus, a
ruler that is one meter long in one frame of reference will not be one meter long in a reference frame that is travelling at a velocity
relative to the first frame. This means that the length of an object is variable depending on the observer.
Units
One of the oldest units of length measurement used in the ancient world was the ‘cubit,’ which was the length of the arm from the
tip of the finger to the elbow. This could then be subdivided into shorter units like the foot, hand (which at 4 inches is still used
today for expressing the height of horses) or finger, or added together to make longer units like the stride. The cubit could vary
considerably due to the different sizes of people.
In the physical sciences and engineering, when one speaks of “units of length”, the word “length” is synonymous with “distance”.
There are several units that are used to measure length. Units of length may be based on lengths of human body parts, the distance
traveled in a number of paces, the distance between landmarks or places on the Earth, or arbitrarily on the length of some fixed
object.In the International System of Units (SI), the basic unit of length is the meter and is now defined in terms of the speed of
light. The centimeter and the kilometer, derived from the meter, are also commonly used units. In U.S. customary units, English or
Imperial system of units, commonly used units of length are the inch, the foot, the yard, and the mile. Units used to denote
distances in the vastness of space, as in astronomy, are much longer than those typically used on Earth and include the astronomical
unit, the light-year, and the parsec.
Length: The metric length of one kilometre is equivalent to the imperial measurement of 0.62137 miles.
Shape
The shape of an object is a description of space that the object takes up; the shape can change if the object is deformed.
learning objectives
Describe effects of deformations, rotations, and magnifications
Shape
The shape of an object located in some space is a geometrical description of the part of that space occupied by the object, as
determined by its external boundary – abstracting from location and orientation in space, size, and other properties such as
color, content, and material composition
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Simple and Complex Shapes
Simple shapes can be described by basic geometry objects such as a set of two or more points, a line, a curve, a plane, a plane
figure (e.g. square or circle), or a solid figure (e.g. cube or sphere). Most shapes occurring in the physical world are complex.
Some, such as plant structures and coastlines, may be so arbitrary as to defy traditional mathematical description – in which case
they may be analyzed by differential geometry, or as fractals.
In geometry, two subsets of a Euclidean space have the same shape if one can be transformed to the other by a combination of
translations, rotations (together also called rigid transformations), and uniform scalings. In other words, the shape of a set of points
is all the geometrical information that is invariant to translations, rotations, and size changes. Having the same shape is an
equivalence relation, and accordingly a precise mathematical definition of the notion of shape can be given as being an equivalence
class of subsets of a Euclidean space having the same shape.
Shapes of physical objects are equal if the subsets of space these objects occupy satisfy the definition above. In particular, the shape
does not depend on the size and placement in space of the object.
Volume
Volume is a measure of the three-dimensional space an object occupies, usually taken in terms of length, width and height.
learning objectives
Explain how is volume measured geometrically
Volume is the quantity of three-dimensional space contained by a closed boundary; it is the space that a substance (solid, liquid, gas
or plasma) or shape occupies or contains. Volume is often quantified numerically using an SI derived unit, the cubic meter.
However, for liquids the unit of volume used is known as the liter (equivalent to 0.001 cubic meters).
Measuring Volume: A measuring cup can be used to measure volumes of liquids. This cup measures volume in units of cups, fluid
ounces and millilitres.
Volume is measured geometrically by multiplying an object’s three dimensions—usually taken as length, width and height. Some
common volumes are taken as follows:
The volume of a cube: length times width times height.
The volume of a cylinder: the cross-sectional area times the height of the cylinder.
The volume of a sphere: 4/3 times the radius cubed times pi.
The volume of a solid can be determined by the volume of liquid it displaces when submerged.
The volume of a container is generally understood as the capacity of the container, meaning the amount of fluid (gas or liquid) the
container can hold, rather than the amount of space the container itself displaces. Measuring cups, as seen in, work by taking a
10.6.2 https://phys.libretexts.org/@go/page/14499
known cross sectional area of a cup and multiplying that by a variable height. Since liquid will always cover the cross section (if
there is enough liquid), adding more liquid will increase the height inside the container.
Liquids take the shape of their container, filling up the minimum height needed. Gases, on the other hand, take up the maximum
amount of volume possible. Thus a measuring cup can accurately measure the volume of a liquid, whereas a gas will always fill the
entire container, more or less uniformly, no matter how little gas there is.
Electric field of positive point charge: The electric field of a positively charged particle points radially away from the charge.
Electric field of negative point charge: The electric field of a negatively charged particle points radially toward the particle.
The reason for these directions can be seen in the derivation of the electric field of a point charge. Let’s first take a look at the
definition of electric field of a point particle:
→ 1 q q
E = ^
r =k ^
r. (10.6.1)
4πϵo r2 2
r
In the above equation, q represents the charge of the particle creating the electric field and the constant k is a result of simply
lumping the constants together. This charge is either positive or negative. If the charge is positive, as shown above, the electric field
will be pointing in a positive radial direction from the charge q (away from the charge) and the following text explains why. The
above equation is defined in radial coordinates which can be seen in.
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Radial Coordinate System: The electric field of a point charge is defined in radial coordinates. The positive r direction points
away from the origin, and the negative r direction points toward the origin. The electric field of a point charge is symmetric with
respect to the θ direction.
If we now place another positive charge, Q (called the test charge), at some radial distance, R, away from the original particle, the
test charge will feel a force given by
→ → 1 q
F = QE = Q ^
r (10.6.2)
4πϵo R2
The thing to keep in mind is that the force above is acting on the test charge Q, in the positive radial direction as defined by the
original charge q. This means that because the charges are both positive and will repel one another, the force on the test charge
points away from the original charge.
If the test charge were negative, the force felt on that charge would be
→ → 1 q
F = Q E = −Q ^
r (10.6.3)
4πϵo R2
Notice that this points in the negative ^r direction, which is toward the original charge. This makes sense because opposite charges
attract and the force on the test charge will tend to push it toward the original positive charge creating the field. The above
mathematical description of the electric field of a point charge is known as Coulomb ‘s Law.
Key Points
Length is typically a measure of the longest dimension of an object.
The deformation of an object is typically a change in length.
The SI unit of length is the meter.
The shape of an object is a representation of the space taken up by the object.
Deformations can change the shape of an object.
Objects that have the same shape can be transformed into each other by rotation or magnification.
Volume is often quantified numerically using an SI derived unit, the cubic meter. However, for liquids the unit of volume used
is known as the liter (equivalent to 0.001 cubic meters).
This can also be understood as the amount of fluid a submerged object displaces.
Volume can be measured for geometrically regular objects by simple formulas. However, more complicated objects are easier to
measure with fluid displacement.
The electric field is a vector field around a charged particle that represents the force that other charged particles would feel
when placed near the particle creating the electric field.
Given a point charge, or a particle of infinitesimal size, that contains a certain charge, electric field lines emanate from equally
in all radial directions.
If the point charge is positive, field lines point away from it; if the charge is negative, field lines point towards it.
10.6.4 https://phys.libretexts.org/@go/page/14499
Key Terms
dimension: A measure of spatial extent in a particular direction, such as height, width or breadth, or depth.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
plane: A level or flat surface.
Euclidean: Adhering to the principles of traditional geometry, in which parallel lines are equidistant.
cross section: A section formed by a plane cutting through an object, usually at right angles to an axis.
dimension: A measure of spatial extent in a particular direction, such as height, width or breadth, or depth.
coulomb’s law: the mathematical equation calculating the electrostatic force vector between two charged particles
vector field: a construction in which each point in a Euclidean space is associated with a vector; a function whose range is a
vector space
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CHAPTER OVERVIEW
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1
CHAPTER OVERVIEW
Front Matter
TitlePage
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11: Fluid Dynamics and Its Applications
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11.1: Overview
learning objectives
Interpret the circulatory system in terms of your knowledge of fluid dynamics
We have discussed many situations in which fluids are static, though there are many situations where fluids flow. For example, a
column of smoke rises from a camp fire, water streams from a fire hose, blood courses through your veins. Why does rising smoke
curl and twist? How does a nozzle increase the speed of water emerging from a hose? How does the body regulate blood flow?
Fluid dynamics, the physics of fluids in motion, allows us to answer these and many other questions.
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A Semipermeable Membrane: A semipermeable membrane with small pores that allow only small molecules to pass through.
learning objectives
Assess the significance of studying volumetric flow in addition to flow velocity
Fluid dynamics is the study of fluids in motion and corresponding phenomena. A fluid in motion has a velocity, just as a solid
object in motion has a velocity. Like the velocity of a solid, the velocity of a fluid is the rate of change of position per unit of time.
In mathematical terms, the velocity of a fluid is the derivative of the position vector of the fluid with respect to time, and is
therefore itself a vector quantity. The flow velocity vector is a function of position, and if the velocity of the fluid is not constant
then it is also a function of time. Equation 1 shows the mathematical expression for the velocity of a fluid in motion. As a vector
quantity, fluid velocity must have at least one non-zero directional component and may have up to three non-zero directional
components. The velocity vector has non-zero components in any orthogonal direction along which motion of the fluid occurs.
Volumetric Flow
In addition to flow velocity, volumetric flow rate is an important quantity in fluid dynamics analysis. Volumetric flow is defined as
the volume of fluid that passes through a given surface per unit time. Qualitatively, Figure 1 shows the notion of volumetric flow
rate regarding a cross-sectional surface of area A. Mathematically, volumetric flow rate is the derivative of the volume of fluid that
passes through a given surface with respect to time; in SI units this is expressed as meterscubed per second. Volumetric flow rate is
related to the flow velocity vector as the surface integral with respect to the surface in question. If the surface area in question is a
flat, plane cross-section, the surface integral reduces as shown in Equation 2, where A is the surface area of the surface in question
and v is the flow velocity of the fluid.
Volumetric Flow Rate: Volumetric Flow Rate Surface Integral and its simplification
Flow Velocity – Volumetric Flow Rate Relation: This figure shows the relation between flow velocity and volumetric flow rate.
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Moreover, only the flow velocity component parallel to the surface normal of the surface in question, or alternatively the flow
velocity component perpendicular to the surface in question contributes to the volumetric flow rate. Figure 1 and Equation 2
illustrate decomposition of the flow velocity vector, making an angle θ with respect to the normal of the surface plane in order to
calculate volumetric flow rate through that surface. Thus, volumetric flow rate for a given fluid velocity and cross-sectional surface
area increases as θ decreases, and is maximized when θ = 0. Volumetric flow rate is an important scalar quantity in fluid dynamics
and is used widely in fluid flow measurements. Volumetric flow rate can be converted to mass flow rate if the density of the fluid is
known. Flow of fluids through a closed system is often analyzed as a hydraulic circuit analogous to electron flow in an electronic
circuit where: 1) the volumetric fluid flow is analogous to the electric current, 2) pressure is analogous to the voltage, and 3) fluid
velocity is analogous to current density.
Key Points
There are many fluids in biology and understanding their behavior in motion is crucial to effective medicine.
The heart pumps a fluid, blood, throughout a series of tubes in the body.
Circulation may be understood through a study of fluid dynamics.
Diffusion is the dominant mechanism by which the exchange of nutrients and waste products occur between the blood and
tissue, and between air and blood in the lungs.
Both osmosis and dialysis are used by the kidneys to cleanse the blood, and the medical application of dialysis through
machinery is important in the treatment of individuals with failing kidney function.
Flow velocity is a vector quantity used to describe the motion of a fluid. It can be easily determined for laminar flow but
complex to determine for turbulent flow.
Volumetric flow rate is the volume of a liquid that passes through a given surface per unit time. It is found from the flow
velocity and the surface area of the surface through which the fluid passes.
Fluid flow through a closed hydraulic system is analyzed much like electron flow through an electronic circuit —where
volumetric flow rate is analogous to current, flow velocity is analogous to current density, and pressure is analogous to voltage
(electrical potential).
Key Terms
vascular: Of, pertaining to, or containing vessels that conduct or circulate fluids (such as blood, lymph, or sap) through the
body of an animal or plant.
osmosis: The net movement of solvent molecules from a region of high solvent potential to a region of lower solvent potential
through a partially permeable membrane.
dialysis: A method of separating molecules or particles of different sizes by differential diffusion through a semipermeable
membrane.
Laminar Flow: Non-turbulent motion of a fluid in which parallel layers have different velocities relative to each other.
Turbulent Flow: The motion of a fluid having local velocities and pressures that fluctuate randomly.
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11.2: Flow in Tubes
learning objectives
Contrast turbulent and laminar flow in constant velocity
Virtually all moving fluids exhibit viscosity, which is a measure of the resistance of a fluid to flow. Viscosity is a basic property
necessary for the analysis of fluid flow.
Representation of Viscosity: A proportionality constant relating an applied shear stress to the resulting shear velocity.
Different fluids exhibit different viscous behavior yet, in this analysis, only Newtonian fluids (fluids with constant velocity
independent of applied shear stress) will be considered. Viscosity in fluids generally decreases with increasing temperature. The
study of the viscous nature of fluids is called rheology.
In analyzing the properties of moving fluids, it is necessary to determine the nature of flow of the fluid. This is generally split into
two categories, laminar and turbulent flow.
Turbulent Flow
Turbulent flow is characterized by irregular flow of a fluid in which there are both inconsistent flow patterns and velocity
variations throughout the volume of the fluid in motion. Analysis of turbulent flow can be very complex and often requires
advanced mathematical analysis to simulate flow in systems on a near case-by-case basis.
It occurs when the Reynolds number is above a certain critical threshold while mixed turbulent–laminar flow occurs within a range
of Reynolds number below this threshold value. At the lower limit of this mixed turbulent–laminar flow Reynolds number region
there is another critical threshold value, below which only laminar flow is possible.
Laminar Flow
Laminar flow consists of a regular-flow pattern with constant-flow velocity throughout the fluid volume and is much easier to
analyze than turbulent flow.
Relative Magnitudes of Velocity Vectors: Laminar fluid flow in a circular pipe at the same direction.
Laminar flow is often encountered in common hydraulic systems, such as where fluid flow is through an enclosed, rigid pipe; the
fluid is incompressible, has constant viscosity, and the Reynolds number is below this lower critical threshold value. It is
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characterized by the flow of a fluid in parallel layers, in which there is no disruption or interaction between the different layers, and
in which each layer flows at a different velocity along the same direction. The variation in velocity between adjacent parallel layers
is due to the viscosity of the fluid and resulting shear forces.
This figure (see ) gives a representation of the relative magnitudes of the velocity vectors of each of these layers for laminar fluid
flow through a circular pipe, in a direction parallel to the pipe axis.
Poiseuille’s Equation: Can be used to determine the pressure drop of a constant viscosity fluid exhibiting laminar flow through a
rigid pipe.
Considering laminar flow of a constant density, incompressible fluid such as for a Newtonian fluid traveling in a pipe, with a
Reynolds number below the upper limit level for fully laminar flow, the pressure difference between two points along the pipe can
be found from the volumetric flow rate, or vice versa. For such a system with a pipe radius of r, fluid viscosity η, distance between
the two points along the pipe Δx = x2 – x1, and the volumetric flow rate Q, of the fluid, the pressure difference between the two
points along the pipe Δp is given by Poiseuille’s equation (see ).
This equation is valid for laminar flow of incompressible fluids only, and may be used to determine a number of properties in the
hydraulic system, if the others are known or can be measured. In practice, Poiseuille’s equation holds for most systems involving
laminar flow of a fluid, except at regions where features disrupting laminar flow, such as at the ends of a pipe, are present.
Poiseuille’s equation as given in this example (see ) is analogous to Ohm ‘s equation for determining the resistance in an electronic
circuit and is of great practical use in hydraulic-circuit analysis.
Blood Flow
Blood flow is the continuous running of blood through the cardiovascular system, which consists of the vessels and the heart.
learning objectives
Outline how normal plasma behaves in a mammalian cardiovascular system
Blood flow is the continuous running of blood through vessels in the cardiovascular system (the mammalian cardiovascular system
is shown in ). Blood is the viscous fluid composed of plasma and cells. The composition of the blood includes plasma, red blood
cells, white blood cells and platelets. In microcirculation, the properties of the blood cells have an important influence on flow.
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An illustrative overview of the mammalian cardiovascular system: Keep in mind that both circular paths are working
simultaneously and not in a sequential manner as the numbering in the illustration might suggest. Both the ventricles are working
together in harmony; as tiny amounts of blood are moving in the pulmonary circuit, the remainder of the blood moves through the
systemic circuit.
The cardiovascular system, which consists of blood vessels and the heart, helps to distribute nutrients, O2, and other products of
metabolism. The blood moves in the blood vessels, while the heart serves as the pump for the blood. The vessel walls of the heart
are elastic and movable, therefore causing the blood and the wall to exert forces on each other and in turn influencing their
respective motion.
The major quantity of interest in describing the motion of blood particles is velocity—the rate of change of the position of an object
with time:
Δx
v = (11.2.1)
Δt
Blood velocities in arteries are higher during systole than during diastole. One parameter to quantify this difference is pulsatility
index (PI), which is equal to the difference between the peak systolic velocity and the minimum diastolic velocity divided by the
mean velocity during the cardiac cycle.
Another important parameter is the acceleration—the rate of change of velocity: a = Δv
Δt
Normal plasma behaves like a Newtonian fluid at rates of shear. Typical values for the viscosity of normal human plasma at 37°C is
1.2Nsm-2. The viscosity of normal plasma varies with temperature in the same way as does that of its solvent, water. (a 5°C
increase of temperature in the physiological range reduces plasma viscosity by about 10%).
The osmotic pressure of the plasma affects the mechanics of the circulation in several ways. An alteration of the osmotic pressure
difference across the membrane of a blood cell causes a shift of water and a change in cell volume. The change, both in shape and
flexibility, affects the mechanical properties of whole blood. Therefore, a change in plasma osmotic pressure alters the hematocrit
(the volume concentration of red cells in the whole blood) by redistributing water between the intravascular and extravascular
spaces. This in turn affects the mechanics of the whole blood.
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Key Points
Viscosity is the resistance of a fluid to flow. Virtually all fluids have viscosity which generally changes as a function of
temperature; although different types of fluids exhibit different types of fluid–shear velocity dependencies.
Laminar flow of a fluid is characterized by its flow in parallel layers in which there is no disruption or interaction between the
different layers, and in which each layer flows at a different velocity along the same direction.
Poiseuille’s equation pertains to moving incompressible fluids exhibiting laminar flow. It relates the difference in pressure at
different spatial points to volumetric flow rate for fluids in motion in certain cases, such as in the flow of fluid through a rigid
pipe.
The major quantity of interest in describing the motion of blood particles is the velocity – the rate of change of the position of
an object with time: v = Δx
Δt
.
Blood velocities in arteries are higher during systole than during diastole.
The mechanics of the circulation depends on osmotic pressure of plasma.
Key Terms
viscosity: The property of a fluid that resists the force which tends to cause it to flow.
shear stress: The external force acting on an object or surface parallel to the slope or plane in which it lies; the stress tending to
produce shear.
vρl
Reynolds Number: A dimensionless number, η
, where v is the fluid velocity, ρ the density, η the viscosity and l a dimension
of the system. The value of the number indicates the type of fluid flow.
systole: The rhythmic contraction of the heart, by which blood is driven through the arteries.
vessel: A tube or canal that carries fluid in an animal or plant.
diastole: The phase or process of relaxation and dilation of the heart chambers, between contractions, during which they fill
with blood; an instance of the process.
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11.3: Bernoulli’s Equation
learning objectives
Adapt Bernoulli’s equation for flows that are either unsteady or compressible
Bernoulli’s Principle: A brief introduction to Bernoulli’s Principle for students studying fluids.
The total mechanical energy of a fluid exists in two forms: potential and kinetic. The kinetic energy of the fluid is stored in static
pressure, psps, and dynamic pressure, 12ρV212ρV2, where \rho is the fluid density in (SI unit: kg/m3) and V is the fluid velocity
(SI unit: m/s). The SI unit of static pressure and dynamic pressure is the pascal.
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Syphoning: Syphoning fluid between two reservoirs. The flow rate out can be determined by drawing a streamline from point ( A )
to point ( C ).
Static pressure is simply the pressure at a given point in the fluid, dynamic pressure is the kinetic energy per unit volume of a fluid
particle. Thus, a fluid will not have dynamic pressure unless it is moving. Therefore, if there is no change in potential energy along
a streamline, Bernoulli’s equation implies that the total energy along that streamline is constant and is a balance between static and
dynamic pressure. Mathematically, the previous statement implies:
1
2
ps + ρV = constant (11.3.1)
2
along a streamline. If changes there are significant changes in height or if the fluid density is high, the change in potential energy
should not be ignored and can be accounted for with,
ΔPE = ρgΔh. (11.3.2)
This simply adds another term to the above version of the Bernoulli equation and results in
1
2
ps + ρV + ρgΔh = constant. (11.3.3)
2
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Torricelli’s Law
Torricelli’s law is theorem about the relation between the exit velocity of a fluid from a hole in a reservoir to the height of fluid
above the hole.
learning objectives
Infer the exit velocity through examining the Bernoulli equation
Torricelli’s law is theorem in fluid dynamics about the relation between the exit velocity of a fluid from a sharp-edged hole in a
reservoir to the height of the fluid above that exit hole. This relationship applies for an “ideal” fluid (inviscid and incompressible)
and results from an exchange of potential energy,
Torricelli’s Principle: A brief introduction to Torricelli’s Principle for students studying fluids.
mghmgh, for kinetic energy,
1
2
2
ρv , at the exit.
This relationship can be derived by applying the Bernoulli equation between the top of the reservoir and the exit hole. Applying
Bernoulli between the top of a reservoir and an exit hole at a height h below the top of the reservoir results in,
Exchange of Energy: Potential energy at the top of the reservoir becomes kinetic energy at the exit.
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1 1
2 2
pt + ρvt + ρg ht = pe + ρve + ρg he (11.3.4)
2 2
where subscript t implies evaluation at the top of the reservoir and subscript e implies evaluation at the exit. If we assume both the
top of reservoir and the exit are open to the atmosphere, the zero for potential energy is at the exit hole, and the fluid velocity at the
top of the reservoir is essentially zero (large reservoir, small hole), we arrive at
1
2
ρg ht = ρve (11.3.5)
2
where again ht is the height difference between the top of the reservoir and the exit hole. Due to the assumption of an ideal fluid, all
forces acting on the fluid are conservative and thus there is an exchange between potential and kinetic energy. The result is that the
velocity acquired by the fluid is the same that a body would acquire when simply dropped from the height ht.
A simple experiment to test Torricelli’s law involves filling a soda bottle with water and puncturing the bottom with a small hole
(about 1 cm in diameter). As the height in the reservoir decreases, the exit velocity will decrease as well. The exit velocity can be
increased by capping the top of the reservoir and pressurizing it.
Toricelli’s Law: The exit velocity depends on the height of the fluid above the exit hole.
Surface Tension
The tendency of the surface of a liquid to resist a force and behave like a membrane and is a result of cohesion between liquid
molecules.
learning objectives
Summarize the cause for different surface tensions at a liquid’s surface
Surface tension is the tendency of a liquid surface to resist forces applied to it. This effect is a result of cohesion of the molecules of
the liquid causing the surface of the liquid to contract to the smallest area possible. This effect is visible in nature with water strider
insects that are able to walk on water. Also, a paper clip or pin can be supported by the surface tension at a water air interface.
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Surface Tension FBD: Force diagrams showing the direction of forces for water supporting a water strider (insect) foot and a pin.
In both cases, the vertical component of the surface tension is enough to support the weight of the object.
In the bulk of the liquid, the molecules are pulled equally in all directions. The molecules at the surface feel a greater attractive
force toward the bulk material than the interface material.
The surface of a liquid is an interface between another fluid, a solid body, or both. Therefore, the surface tension will be a property
of the interface rather than simply the liquid. Adhesion describes the attractive force between molecules of different types. The
surface of a liquid in a container is an interface between the liquid, the air, and the container. Where the surfaces meet, forces must
be in equilibrium. This results in a contact angle at the interface. The contact angle is measured in the liquid and depends on the
relative strength of cohesive forces in the liquid and adhesive forces between the liquid and interface materials. If liquid molecules
are strongly attracted to the molecules of the solid surface (adhesive forces > cohesive forces), the drop will tend to spread out and
the contact angle will be close to zero degrees. If the cohesive forces are greater than the adhesive forces, the resulting contact
angles will be large and will form a more circular drop.
Water Droplet on Leaf: When a water droplet forms on a leaf, the cohesive forces between the water molecules are greater than
the adhesive forces between the water and leaf surface. The leaf is a hydrophobic surface.
Contact Angle: The contact angle is the angle, measured in the fluid, that results when a liquid-gas interface, meets a solid
surface.
When the liquid is water, a surface where the contact angle is small is said to be hydrophilic. Large contact angles are present on
hydrophobic surfaces. The contact angle determines the wettability of the surface.
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Key Points
The simplest form of Bernoulli’s equation (steady and incompressible flow) states that the sum of mechanical energy, potential
energy and kinetic energy, along a streamline is constant. Therefore, any increase in one form results in a decrease in the other.
Bernoulli’s equation considers only pressure and gravitational forces acting on the fluid particles. Therefore, if there is no
change in height along a streamline, Bernoulli’s equation becomes a balance between static pressure and velocity.
The steady-state, incompressible Bernoulli equation, can be derived by integrating Newton’s 2nd law along a streamline.
Torricelli’s law applies to an inviscid, incompressible fluid (“ideal” fluid).
You can ascertain results from applying the Bernoulli equation between the top of the reservoir and the exit hole.
The relationship arises from an exchange of potential energy at the top of the reservoir to kinetic energy at the exit.
The final kinetic energy is equivalent to what a solid body would acquire when falling from height h.
Surface tension is a result of cohesion between the molecules of the liquid. The molecules at the surface of the liquid feel an
attractive force pulling them toward the bulk of the liquid more than the solid or fluid at the interface.
When a liquid-solid-gas interface is encountered, the contact angle represents a measure of the relative strength of adhesive and
cohesive forces.
The contact angle determines the wettability of a surface.
Key Terms
viscosity: A quantity expressing the magnitude of internal friction in a fluid, as measured by the force per unit area resisting
uniform flow.
Ideal Fluid: An inviscid and incompressible fluid
incompressible: Unable to be compressed or condensed.
inviscid: A fluid with zero viscosity (internal friction). In reality viscosity is always present. However, it is often very small
compared with other forces (e.g. gravity, pressure) and for common fluids (water and air) the fluid can be approximated as
having zero viscosity.
cohesion: Various intermolecular forces that hold solids and liquids together.
wettability: The ability of a solid surface to reduce the surface tension of a liquid in contact with it such that it spreads over the
surface and wets it.
adhesion: The ability of a substance to stick to an unlike substance.
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11.4: Other Applications
learning objectives
Predict if flow will be laminar or turbulent
It is possible to predict if flow will be laminar or turbulent. At low velocity, flow in a very smooth tube or around a smooth,
streamlined object will be laminar. At high velocity, even the flow in a smooth tube or around a smooth object will experience
turbulence. However, between low and high velocity, flow is more difficult to predict. In fact, at intermediate velocities, flow may
oscillate back and forth indefinitely between laminar and turbulent.
An occlusion (narrowing) of an artery, such as shown in, is likely to cause turbulence because of the irregularity of the blockage, as
well as the complexity of blood as a fluid. Turbulence in the circulatory system (such as aneurysms, or ballooning of arteries) is
noisy and can sometimes be detected with a stethoscope (such as when measuring diastolic pressure in the upper arm’s partially
collapsed brachial artery). These turbulent sounds, at the onset of blood flow when the cuff pressure becomes sufficiently small, are
called Korotkoff sounds. Heart murmurs, consistent with their name, are sounds produced by turbulent flow around damaged and
insufficiently closed heart valves. Another method of detecting this type of turbulence is ultrasound, used as a medical indicator in
a process analogous to Doppler-shift radar (used to detect storms).
Turbulent Flow in an Artery: Flow is laminar in the large part of this blood vessel and turbulent in the part narrowed by plaque,
where velocity is high. In the transition region, the flow can oscillate chaotically between laminar and turbulent flow.
Turbulence manifests in other areas, with varying causes. During an airplane flight, for example, the turbulence experienced is due
to the mixing of warm and cold air in the atmosphere, causing the airplane to shake. The mixing currents in oceans creates a similar
effect.
The phenomenon of turbulent air flow must be accounted for in many applications. For example, race cars are unable to follow
each other around fast corners because the leading car creates turbulent air flow in its wake (this can lead to under-steering).
Industrial equipment, such as pipes, ducts, and heat exchangers are often designed to induce the flow regime of interest (laminar or
turbulent). When flow is turbulent, particles exhibit additional transverse motion. This enhances the rate of energy and momentum
exchange between them, increasing the heat transfer. Turbulent flow is thus desirable in applications where a relatively cool fluid is
mixed with a warmer fluid to reduce the temperature of the warmer fluid.
It is imperative to take into account turbulent flow when designing certain structures, such as a bridge support, as shown in. In the
late summer and fall, when river flow is slow, water flows smoothly around the support legs. In the spring, when the flow is faster,
the flow may start off laminar but it is quickly separated from the leg and becomes turbulent. The bridge supports must be designed
so that they can withstand the turbulent flow of the water in the spring.
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Longtown Bridge: Turbulent flow is visible around the bridge supports of the Longtown bridge.
learning objectives
Assess the relationship of the parameters to one another in determining the inertia of an object moving in fluid
Overview
A moving object in a viscous fluid is equivalent to a stationary object in a flowing fluid stream. (For example, when you ride a
bicycle at 10 m/s in still air, you feel the air in your face exactly as if you were stationary in a 10-m/s wind. ) Flow of the stationary
fluid around a moving object may be laminar, turbulent, or a combination of the two. Just as with flow in tubes, it is possible to
predict when a moving object creates turbulence. We use another form of the Reynolds number N′R, defined for an object moving
in a fluid to be
ρvL
N‘R = (11.4.1)
η
where L is a characteristic length of the object (a sphere’s diameter, for example), the fluid density, its viscosity, and v the object’s
speed in the fluid. If N′R is less than about 1, flow around the object can be laminar, particularly if the object has a smooth shape.
The transition to turbulent flow occurs for N′R between 1 and about 10, depending on surface roughness and so on. Depending on
the surface, there can be a turbulent wake behind the object with some laminar flow over its surface. For an N′R between 10 and
10^6, the flow may be either laminar or turbulent and may oscillate between the two. For N′R greater than about 10^6, the flow is
entirely turbulent, even at the surface of the object. (See. ) Laminar flow occurs mostly when the objects in the fluid are small, such
as raindrops, pollen, and blood cells in plasma.
Motion of an object in a viscous fluid.: (a) Motion of this sphere to the right is equivalent to fluid flow to the left. Here the flow is
laminar with N′R less than 1. There is a force, called viscous drag FV, to the left on the ball due to the fluid’s viscosity. (b) At a
higher speed, the flow becomes partially turbulent, creating a wake starting where the flow lines separate from the surface. Pressure
in the wake is less than in front of the sphere, because fluid speed is less, creating a net force to the left F′V that is significantly
greater than for laminar flow. Here N′R is greater than 10. (c) At much higher speeds, where N′R is greater than 10^6, flow
becomes turbulent everywhere on the surface and behind the sphere. Drag increases dramatically.
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Viscous Drag
One of the consequences of viscosity is a resistance force called viscous drag FVFV that is exerted on a moving object. This force
typically depends on the object’s speed (in contrast with simple friction). Experiments have shown that for laminar flow (N′R less
than about one) viscous drag is proportional to speed, whereas for N′R between about 10 and 106, viscous drag is proportional to
speed squared. (This relationship is a strong dependence and is pertinent to bicycle racing, where even a small headwind causes
significantly increased drag on the racer. Cyclists take turns being the leader in the pack for this reason. ) For N′R greater than 106,
drag increases dramatically and behaves with greater complexity. For laminar flow around a sphere, FVFV is proportional to fluid
viscosity, the object’s characteristic size L, and its speed v. All of which makes sense—the more viscous the fluid and the larger the
object, the more drag we expect. Recall Stoke’s law FS=6πrηvFS=6πrηv. For the special case of a small sphere of radius R, moving
slowly in a fluid of viscosity, the drag force FSFS is given by
FS=6πRηvFS=6πRηv.
learning objectives
Predict the role diffusion plays in blood transport throughout the body
Diffusion
Atoms and molecules are in constant motion at any temperature. In fluids they move about randomly even in the absence of
macroscopic flow.
Diffusion is the movement of substances due to random thermal molecular motion. Fluids, like fish fumes or odors entering ice
cubes, can even diffuse through solids. Diffusion is a slow process over macroscopic distances. The densities of common materials
are great enough that molecules cannot travel very far before having a collision that can scatter them in any direction, including
straight backward. More massive molecules diffuse more slowly.
Another interesting point is that the diffusion rate for oxygen in air is much greater than for oxygen in water. In water, an oxygen
molecule makes many more collisions in its random walk and is slowed considerably. In water, an oxygen molecule moves only
about 40μm in 1 s. (Each molecule actually collides about 1010 times per second!). Finally, note that diffusion constants increase
with temperature, because average molecular speed increases with temperature. This is because the average kinetic energy of
molecules, 1/2mv2, is proportional to absolute temperature. Because diffusion is typically very slow, its most important effects
occur over small distances. For example, the cornea of the eye gets most of its oxygen by diffusion through the thin tear layer
covering it.
If you very carefully place a drop of food coloring in a still glass of water, it will slowly diffuse into the colorless surroundings
until its concentration is the same everywhere. This type of diffusion is called free diffusion, because there are no barriers
inhibiting it. Let us examine its direction and rate. Molecular motion is random in direction, and so simple chance dictates that
more molecules will move out of a region of high concentration than into it. The net rate of diffusion is higher initially than after
the process is partially completed. The rate of diffusion is proportional to the concentration difference. Many more molecules will
leave a region of high concentration than will enter it from a region of low concentration. In fact, if the concentrations were the
same, there would be no net movement. The rate of diffusion is also proportional to the diffusion constant D, which is determined
experimentally. Many of the factors that affect the rate are hidden in the diffusion constant D. For example, temperature and
cohesive and adhesive forces all affect values of D. Diffusion is the dominant mechanism by which the exchange of nutrients and
waste products occur between the blood and tissue, and between air and blood in the lungs. In the evolutionary process, as
organisms became larger, they needed quicker methods of transportation than net diffusion, because of the larger distances involved
in the transport, leading to the development of circulatory systems. Less sophisticated, single-celled organisms still rely totally on
diffusion for the removal of waste products and the uptake of nutrients.
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Food Coloring: Food coloring spreading on a thin water film.
Diffusion: (a) Two sugar-water solutions of different concentrations, separated by a semipermeable membrane that passes water
but not sugar. Osmosis will be to the right, since water is less concentrated there. (b) The fluid level rises until the back pressure
ρgh equals the relative osmotic pressure; then, the net transfer of water is zero.
learning objectives
Contrast systole and diastole in cardiovascular circulation
11.4.4 https://phys.libretexts.org/@go/page/14504
located at the bottom of the aorta and pulmonary artery. Strong tendinous chords attached to valves prevent them from turning
inside out when they close.
The human heart will undergo over 3 billion contraction cycles during a normal lifetime. A complete cardiac cycle is one round of
the heart pumping blood and consists of two parts: systole (contraction of the heart muscle) and diastole (relaxation of the heart
muscle). During the cycle, the top half of the heart works as one unit, while the bottom half of the heart works as one unit.
The heart beat can be heard as a sound that the valves make when they close. The ‘lub’ sound is made when the atrio ventricular
valves close and the ‘dub’ sound is made when the semi lunar valves close. Blood pressure is produced by the left ventricle
contractions.The rhythm of ventricle diastole, often just referred to as diastole, causes the pulse, which can be felt by holding two
fingers to the side of the throat.
Key Points
For low velocity, flow in a smooth tube will be laminar.
At higher velocities or if there are obstructions, the flow turns turbulent.
Turbulent flow is very chaotic, with rapid variations in velocity and pressure.
Viscous fluids exert a resitive force on objects attempting to move through them.
This resistive force is called viscous drag and is proportional to the viscosity of the fluid and the motion of the object.
An object moving in a fluid can be thought of as a stationary object in a moving fluid.
Diffusion is the movement of molecules due to random thermal motion.
Osmosis is the movement of molecules due to different concentrations. Molecules will move from regions of high
concentrations to lower concentrations.
These transport phenomena can take place through membranes if the pressure is great enough.
Key Terms
turbulent: Being in, or causing, disturbance or unrest.
streamlined: Designed to offer little resistance to the flow of fluid, especially by having sleek, graceful lines.
laminar: Of fluid motion, smooth and regular, flowing as though in different layers.
viscosity: A quantity expressing the magnitude of internal friction in a fluid, as measured by the force per unit area resisting
uniform flow.
turbulence: Disturbance in a gas or fluid, characterized by evidence of internal motion or unrest.
diffusion: the intermingling of the molecules of a fluid due to random thermal agitation
ventricle: One of two lower chambers of the heart.
contraction: A reversible reduction in size.
atrium: An upper chamber of the heart that receives blood from the veins and forces it into a ventricle. In higher vertebrates,
the right atrium receives blood from the superior vena cava and inferior vena cava, and the left atrium receives blood from the
left and right pulmonary veins.
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CHAPTER OVERVIEW
Back Matter
Index
1 https://phys.libretexts.org/@go/page/26949
Index
A angular frequency B
aberration 15.3: Periodic Motion ballistics
24.3: Lenses Angular momentum 5.4: Types of Forces in Nature
25.2: Other Optical Instruments 29.3: Atomic Physics and Quantum Mechanics basal metabolic rate
absolute space 5.9: Angular vs. Linear Quantities
9.6: Conservation of Angular Momentum 6.6: Power
27.1: Introduction 9.7: Vector Nature of Rotational Kinematics base
absolute zero angular motion 17.6: Applications of Electrostatics
12.2: Temperature and Temperature Scales 8.1: Introduction battery
12.4: Ideal Gas Law
14.4: Entropy angular position 19.2: Electric Current
9.1: Quantities of Rotational Kinematics 20.1: Overview
14.5: The Third Law of Thermodynamics
acceleration angular velocity Bell's theorem
5.3: Velocity, Acceleration, and Force 28.1: History and Quantum Mechanical Quantities
1.2: Units
2.3: Acceleration 5.6: Kepler’s Laws Bernoulli’s equation
2.5: Free-Falling Objects 5.9: Angular vs. Linear Quantities 11.3: Bernoulli’s Equation
3.2: Vectors 8.1: Introduction
9.1: Quantities of Rotational Kinematics
Beta decay
4.3: Newton’s Laws
5.3: Velocity, Acceleration, and Force 9.10: Conservation of Energy 30.2: Radioactivity
achromatic 9.3: Rotational Kinematics bilateral symmetry
9.5: Rotational Kinetic Energy 3.3: Projectile Motion
25.2: Other Optical Instruments 9.7: Vector Nature of Rotational Kinematics
achromatic doublet binary data
annihilation 26.4: Applications of Wave Optics
24.3: Lenses 27.3: Relativistic Quantities
action potential black body
antimatter 13.4: Methods of Heat Transfer
19.6: Electricity in the World 27.3: Relativistic Quantities 28.1: History and Quantum Mechanical Quantities
adhesion antinode 29.1: Overview
11.3: Bernoulli’s Equation 16.1: Introduction 29.2: The Early Atom
Adiabatic 16.5: Further Topics Black body radiation
12.4: Ideal Gas Law aperture 28.1: History and Quantum Mechanical Quantities
adiabatic index 25.1: The Human Eye blackbody
13.2: Specific Heat 25.2: Other Optical Instruments 13.5: Global Warming
aether application Boltzmann constant
27.2: Consequences of Special Relativity 1.1: The Basics of Physics 12.5: Kinetic Theory
afocal system approximation boundary condition
24.3: Lenses 1.3: Significant Figures and Order of Magnitude 15.6: Wave Behavior and Interaction
alpha particle Archimedes' Principle Boyle's Law
29.2: The Early Atom 10.3: Archimedes’ Principle 14.2: The First Law of Thermodynamics
30.2: Radioactivity artificial satellite brachytherapy
alternating current 5.6: Kepler’s Laws 30.4: Applications of Nuclear Physics
20.5: RC Circuits asteroid breakdown
Alveoli 14.4: Entropy 18.4: Capacitors and Dielectrics
10.2: Density and Pressure 5.1: Introduction to UCM and Gravitation
brine
AM Radio Waves astronomical unit 12.2: Temperature and Temperature Scales
23.1: The Electromagnetic Spectrum 5.6: Kepler’s Laws
Brownian motion
ampere atom 12.1: Introduction
19.2: Electric Current 12.1: Introduction 12.4: Ideal Gas Law
21.5: Magnetic Fields, Magnetic Forces, and atomic number buoyancy
Conductors 29.5: Multielectron Atoms 10.2: Density and Pressure
Ampere's law atomic orbital buoyant force
21.1: Magnetism and Magnetic Fields 28.2: Applications of Quantum Mechanics 10.3: Archimedes’ Principle
amplifier 29.5: Multielectron Atoms
22.3: Applications of Induction and EM Waves atomic spectra C
amplitude 30.1: The Nucleus
15.3: Periodic Motion atrium calorie
15.7: Waves on Strings 14.2: The First Law of Thermodynamics
11.4: Other Applications
16.2: Sound Intensity and Level
Avogadro’s number calorimeter
16.4: Interactions with Sound Waves 13.2: Specific Heat
Aneroid Barometer 12.4: Ideal Gas Law
29.2: The Early Atom calorimetry
10.2: Density and Pressure 13.1: Introduction
axis
angular Capacitance
3.2: Vectors
9.2: Angular Acceleration 18.4: Capacitors and Dielectrics
axon
angular acceleration
18.2: Equipotential Surfaces and Lines
8.8: Torque and Angular Acceleration 18.3: Point Charge
9.1: Quantities of Rotational Kinematics
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capacitor component Coulomb force
17.4: The Electric Field Revisited 3.2: Vectors 21.3: Magnetic Force on a Moving Electric Charge
18.4: Capacitors and Dielectrics compression 6.5: Potential Energy and Conservation of Energy
20.3: Kirchhoff’s Rules Coulomb’s Law
15.5: Waves
20.5: RC Circuits
compressive stress 10.6: Deformation of Solids
Carnot cycle 17.4: The Electric Field Revisited
8.5: Applications of Statics
14.4: Entropy Coulomb's Law
cathode computed tomography
17.3: Coulomb’s Law
30.4: Applications of Nuclear Physics
29.2: The Early Atom crest
CCD concave lens
15.5: Waves
24.3: Lenses
29.4: Applications of Atomic Physics Critical temperature
cell membrane concave mirror
19.3: Resistance and Resistors
24.4: Mirrors
18.2: Equipotential Surfaces and Lines critically damped
18.3: Point Charge concentration
15.4: Damped and Driven Oscillations
center of mass 12.9: Diffusion
cryptography
5.5: Newton’s Law of Universal Gravitation conductive medium
28.1: History and Quantum Mechanical Quantities
7.5: Center of Mass 19.1: Overview
8.3: Stability conductor crystallography
8.7: The Center of Gravity 25.2: Other Optical Instruments
17.1: Overview
centrifugal force 17.4: The Electric Field Revisited
29.2: The Early Atom
5.4: Types of Forces in Nature 18.2: Equipotential Surfaces and Lines Curie temperature
centripetal 18.3: Point Charge 21.2: Magnets
18.4: Capacitors and Dielectrics Current
29.2: The Early Atom
23.3: Applications of EM Waves
5.1: Introduction to UCM and Gravitation 19.2: Electric Current
5.2: Non-Uniform Circular Motion conservation 20.3: Kirchhoff’s Rules
5.3: Velocity, Acceleration, and Force 6.5: Potential Energy and Conservation of Energy 21.5: Magnetic Fields, Magnetic Forces, and
centripetal acceleration 7.1: Introduction Conductors
15.3: Periodic Motion Conservation of Energy Cyclotron
CFC 7.3: Collisions 21.4: Motion of a Charged Particle in a Magnetic
conservation of momentum principle Field
14.3: The Second Law of Thermodynamics
characteristic time constant 7.3: Collisions cyclotron frequency
conservative force 21.4: Motion of a Charged Particle in a Magnetic
22.2: AC Circuits
Field
chemical reaction 6.5: Potential Energy and Conservation of Energy
12.1: Introduction constant velocity
D
Chromatic aberration 3.1: Motion in Two Dimensions
25.2: Other Optical Instruments constructive interference damping
16.5: Further Topics
circuit 15.6: Wave Behavior and Interaction
18.4: Capacitors and Dielectrics
16.5: Further Topics de Broglie wavelength
25.2: Other Optical Instruments 29.4: Applications of Atomic Physics
Circular Motion 26.2: Diffraction
5.3: Velocity, Acceleration, and Force
Decay
continuity 30.2: Radioactivity
cladding 10.5: Fluids in Motion
24.2: Reflection, Refraction, and Dispersion
decibels (units)
contraction 16.2: Sound Intensity and Level
classical electrodynamics 11.4: Other Applications
23.2: Electromagnetic Waves and their Properties
deformation
contrast 15.2: Hooke’s Law
Classical mechanics 25.2: Other Optical Instruments 8.6: Elasticity, Stress, Strain, and Fracture
27.3: Relativistic Quantities 26.4: Applications of Wave Optics
degeneracy
classical physics control rod 14.5: The Third Law of Thermodynamics
16.3: Doppler Effect and Sonic Booms 30.4: Applications of Nuclear Physics
degree of freedom
closed system conversion 12.1: Introduction
7.1: Introduction 1.2: Units 7.3: Collisions
cochlea convex demagnetization
16.2: Sound Intensity and Level 25.2: Other Optical Instruments 14.5: The Third Law of Thermodynamics
coherence convex lens destructive interference
28.2: Applications of Quantum Mechanics 24.3: Lenses 15.6: Wave Behavior and Interaction
coherent convex mirror 16.5: Further Topics
16.4: Interactions with Sound Waves 24.4: Mirrors 25.2: Other Optical Instruments
26.1: Superposition and Interference 26.2: Diffraction
coordinate
cohesion diagram
3.2: Vectors
11.3: Bernoulli’s Equation 2.4: Problem-Solving for Basic Kinematics
Coordinate axes
combination circuit dialysis
3.2: Vectors
20.2: Resistors in Series and Parallel 11.1: Overview
corona discharge
combustion diamagnetism
17.6: Applications of Electrostatics
13.2: Specific Heat 21.6: Applications of Magnetism
coulomb
complex numbers diastole
17.1: Overview
19.5: Alternating Currents 18.1: Overview 11.2: Flow in Tubes
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dielectric drift velocity electromotive force (emf)
17.1: Overview 19.2: Electric Current 20.3: Kirchhoff’s Rules
17.2: Shelding and Charging Through Induction 21.5: Magnetic Fields, Magnetic Forces, and 22.1: Magnetic Flux, Induction, and Faraday’s Law
17.4: The Electric Field Revisited Conductors electron affinity
18.4: Capacitors and Dielectrics dynamic contrast ratio 28.2: Applications of Quantum Mechanics
24.2: Reflection, Refraction, and Dispersion
25.1: The Human Eye electron gun
differential dynamics 18.5: Applications
12.8: Thermal Stresses
1.4: Solving Physics Problems electron shell
differential equation 4.7: Further Applications of Newton’s Laws
21.6: Applications of Magnetism
20.5: RC Circuits dynamo 29.5: Multielectron Atoms
23.2: Electromagnetic Waves and their Properties
21.1: Magnetism and Magnetic Fields electron volt
diffraction
18.1: Overview
25.2: Other Optical Instruments
26.2: Diffraction
E electrostatic force
28.1: History and Quantum Mechanical Quantities eardrum 17.3: Coulomb’s Law
29.2: The Early Atom 16.2: Sound Intensity and Level element
diffusion eccentricity 29.5: Multielectron Atoms
11.4: Other Applications 5.6: Kepler’s Laws elementary charge
12.9: Diffusion
19.6: Electricity in the World
elastic 21.5: Magnetic Fields, Magnetic Forces, and
7.2: Conservation of Momentum Conductors
dimension 8.6: Elasticity, Stress, Strain, and Fracture elementary particle
1.4: Solving Physics Problems
10.6: Deformation of Solids
Elastic Collision 28.1: History and Quantum Mechanical Quantities
7.3: Collisions 7.1: Introduction emissivity
7.3: Collisions
diode 13.4: Methods of Heat Transfer
30.2: Radioactivity
elastic potential energy 23.1: The Electromagnetic Spectrum
30.4: Applications of Nuclear Physics 15.3: Periodic Motion Energy
diopter elasticity 15.5: Waves
25.2: Other Optical Instruments 15.2: Hooke’s Law 6.1: Introduction
dipole 16.1: Introduction Enthalpy
24.2: Reflection, Refraction, and Dispersion
electric charge 10.1: Introduction
17.1: Overview 13.2: Specific Heat
dipole moment 17.5: Electric Flux and Gauss’s Law entropy
17.1: Overview 19.1: Overview
18.1: Overview 12.2: Temperature and Temperature Scales
20.3: Kirchhoff’s Rules 14.3: The Second Law of Thermodynamics
21.6: Applications of Magnetism 21.6: Applications of Magnetism 14.4: Entropy
direct current electric current epistemological
20.1: Overview 19.2: Electric Current
20.5: RC Circuits 28.1: History and Quantum Mechanical Quantities
electric displacement field epithelium
direction of propagation 17.5: Electric Flux and Gauss’s Law
15.5: Waves 16.4: Interactions with Sound Waves
Electric field equation
discharge 17.1: Overview
17.1: Overview 4.5: Problem-Solving
17.5: Electric Flux and Gauss’s Law
17.2: Shelding and Charging Through Induction 18.1: Overview equilibrium
disorder 21.5: Magnetic Fields, Magnetic Forces, and 10.2: Density and Pressure
14.4: Entropy Conductors 12.6: Phase Changes
22.4: Magnetic Fields and Maxwell Revisited 12.9: Diffusion
Dispersion 13.6: Phase Equilbrium
26.3: Further Topics
electric potential
15.4: Damped and Driven Oscillations
18.1: Overview
displacement 17.4: The Electric Field Revisited
16.4: Interactions with Sound Waves
electrical circuit 4.7: Further Applications of Newton’s Laws
20.1: Overview 8.2: Conditions for Equilibrium
2.1: Basics of Kinematics
3.2: Vectors electrical resistance equipotential
dissipative force 19.1: Overview 18.1: Overview
18.2: Equipotential Surfaces and Lines
15.3: Periodic Motion electrocardiogram 18.3: Point Charge
6.5: Potential Energy and Conservation of Energy 19.6: Electricity in the World
distortion escape velocity
electromagnet 5.8: Energy Conservation
25.2: Other Optical Instruments 21.1: Magnetism and Magnetic Fields
diurnal Euclidean
electromagnetic force 10.6: Deformation of Solids
5.4: Types of Forces in Nature 29.2: The Early Atom
Doppler effect evaporation
electromagnetic radiation 13.6: Phase Equilbrium
16.3: Doppler Effect and Sonic Booms 26.3: Further Topics
23.2: Electromagnetic Waves and their Properties 30.2: Radioactivity
exponent
Doppler shift 6.8: Further Topics 1.3: Significant Figures and Order of Magnitude
29.2: The Early Atom Electromagnetic waves external beam therapy
dosimeter 23.2: Electromagnetic Waves and their Properties 30.4: Applications of Nuclear Physics
30.4: Applications of Nuclear Physics electromotive force
dot product 20.1: Overview
6.2: Work Done by a Constant Force 20.2: Resistors in Series and Parallel
22.4: Magnetic Fields and Maxwell Revisited
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F frequency ground state
Faraday constant 15.3: Periodic Motion 28.2: Applications of Quantum Mechanics
15.5: Waves group velocity
12.4: Ideal Gas Law 16.1: Introduction
Faraday shield 15.5: Waves
16.3: Doppler Effect and Sonic Booms
17.6: Applications of Electrostatics 16.4: Interactions with Sound Waves gyroradius
22.3: Applications of Induction and EM Waves 21.4: Motion of a Charged Particle in a Magnetic
Faraday's law of induction 23.2: Electromagnetic Waves and their Properties Field
20.1: Overview
22.1: Magnetic Flux, Induction, and Faraday’s Law
friction
ferromagnet 4.3: Newton’s Laws H
4.7: Further Applications of Newton’s Laws
22.1: Magnetic Flux, Induction, and Faraday’s Law 6.5: Potential Energy and Conservation of Energy
harmonic oscillator
ferromagnetic 7.3: Collisions 28.1: History and Quantum Mechanical Quantities
21.1: Magnetism and Magnetic Fields frigorific mixture heat
21.6: Applications of Magnetism 12.2: Temperature and Temperature Scales 14.1: Introduction
fibrillation fundamental thermodynamic relation 14.2: The First Law of Thermodynamics
19.6: Electricity in the World 13.2: Specific Heat
Heat capacity
fictitious force fusion 13.2: Specific Heat
5.4: Types of Forces in Nature 30.3: Quantum Tunneling and Conservation Laws
heat engine
field of view 30.4: Applications of Nuclear Physics 14.3: The Second Law of Thermodynamics
6.8: Further Topics 14.4: Entropy
25.1: The Human Eye
first law of thermodynamics heat of reaction
12.4: Ideal Gas Law G 13.2: Specific Heat
14.3: The Second Law of Thermodynamics Galilean transformation heat pump
6.8: Further Topics 14.4: Entropy
27.3: Relativistic Quantities
fission galvanometer heat transfer
30.3: Quantum Tunneling and Conservation Laws 13.1: Introduction
20.4: Voltmeters and Ammeters
30.4: Applications of Nuclear Physics helical
22.1: Magnetic Flux, Induction, and Faraday’s Law
6.8: Further Topics
fluid gamma (γ) rays 26.4: Applications of Wave Optics
23.1: The Electromagnetic Spectrum helical motion
4.7: Further Applications of Newton’s Laws
29.2: The Early Atom 21.4: Motion of a Charged Particle in a Magnetic
fluidity 30.2: Radioactivity Field
10.1: Introduction 30.4: Applications of Nuclear Physics
hertz
flux Gamma decay 16.1: Introduction
22.1: Magnetic Flux, Induction, and Faraday’s Law 23.1: The Electromagnetic Spectrum
Hohmann transfer orbit
23.2: Electromagnetic Waves and their Properties gamma ray 5.6: Kepler’s Laws
flux density 29.3: Atomic Physics and Quantum Mechanics
Hooke's Law
18.3: Point Charge gas constant 6.5: Potential Energy and Conservation of Energy
flux intensity 12.4: Ideal Gas Law
humidity
18.2: Equipotential Surfaces and Lines Gauge Pressure 12.6: Phase Changes
FM 10.2: Density and Pressure
Huygen's principle
23.1: The Electromagnetic Spectrum general relativity 16.5: Further Topics
focal point 27.4: Implications of Special Relativity
hydraulic press
24.3: Lenses geometric optics 10.2: Density and Pressure
Force 24.1: Overview
hydrogen bond
15.4: Damped and Driven Oscillations geothermal 12.3: Thermal Expansion
4.1: Introduction
14.4: Entropy
4.2: Force and Mass
6.3: Work Done by a Variable Force gimbal I
7.3: Collisions 9.7: Vector Nature of Rotational Kinematics
8.2: Conditions for Equilibrium ideal energy
gluon
fossil fuels 12.4: Ideal Gas Law
30.1: The Nucleus
6.7: CASE STUDY: World Energy Use ideal fluid
gradient
Fourier analysis 11.3: Bernoulli’s Equation
5.4: Types of Forces in Nature
15.1: Introduction Ideal gas
grating
frame of reference 10.2: Density and Pressure
28.1: History and Quantum Mechanical Quantities 12.1: Introduction
2.1: Basics of Kinematics Gravitational acceleration 12.2: Temperature and Temperature Scales
22.1: Magnetic Flux, Induction, and Faraday’s Law 12.4: Ideal Gas Law
4.4: Other Examples of Forces
free electron laser Gravitational Force 12.5: Kinetic Theory
28.2: Applications of Quantum Mechanics 14.2: The First Law of Thermodynamics
5.5: Newton’s Law of Universal Gravitation
free fall image distance
gravity 24.3: Lenses
2.5: Free-Falling Objects
17.1: Overview
3.3: Projectile Motion
impedance
5.7: Gravitational Potential Energy 20.5: RC Circuits
22.2: AC Circuits
greenhouse effect
13.5: Global Warming
impulse
7.1: Introduction
greenhouse gas
13.5: Global Warming
incandescence
26.4: Applications of Wave Optics
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incident ray isentropic lens
26.3: Further Topics 14.5: The Third Law of Thermodynamics 25.2: Other Optical Instruments
incline isolated system 26.1: Superposition and Interference
4.7: Further Applications of Newton’s Laws 13.1: Introduction Lenz's law
incompressible 6.5: Potential Energy and Conservation of Energy 22.2: AC Circuits
10.2: Density and Pressure isotope leverage
10.5: Fluids in Motion 30.4: Applications of Nuclear Physics 8.5: Applications of Statics
11.3: Bernoulli’s Equation isotropic line element
index of refraction 12.3: Thermal Expansion 27.4: Implications of Special Relativity
24.2: Reflection, Refraction, and Dispersion linear momentum
26.3: Further Topics
J 7.1: Introduction
induction linear thermal expansion coefficient
22.1: Magnetic Flux, Induction, and Faraday’s Law
joint
8.5: Applications of Statics 12.3: Thermal Expansion
5.5: Newton’s Law of Universal Gravitation
inductor linear velocity
17.2: Shelding and Charging Through Induction K 9.3: Rotational Kinematics
22.1: Magnetic Flux, Induction, and Faraday’s Law kelvin Longitudinal
22.2: AC Circuits 15.5: Waves
16.1: Introduction
22.4: Magnetic Fields and Maxwell Revisited Lorentz factor
inelastic Kelvin scale
12.2: Temperature and Temperature Scales 27.2: Consequences of Special Relativity
7.1: Introduction 27.3: Relativistic Quantities
7.2: Conservation of Momentum kilocalorie
Lorentz force
inelastic scattering 13.1: Introduction
17.3: Coulomb’s Law
29.2: The Early Atom Kinematic 22.1: Magnetic Flux, Induction, and Faraday’s Law
inertia 2.3: Acceleration Lorentz transformations
2.4: Problem-Solving for Basic Kinematics
1.2: Units 27.1: Introduction
3.1: Motion in Two Dimensions
4.3: Newton’s Laws 27.3: Relativistic Quantities
9.2: Angular Acceleration
9.5: Rotational Kinetic Energy 27.4: Implications of Special Relativity
9.8: Problem Solving
inertial frame kinematics Lorenz Invariance
5.4: Types of Forces in Nature 27.1: Introduction
2.1: Basics of Kinematics
instantaneous Kinetic Energy luminance
2.2: Speed and Velocity 25.1: The Human Eye
10.2: Density and Pressure
insulator 15.2: Hooke’s Law luminiferous aether
17.1: Overview 18.1: Overview 27.1: Introduction
interfere 4.7: Further Applications of Newton’s Laws
5.8: Energy Conservation
16.4: Interactions with Sound Waves 6.8: Further Topics M
interference 7.3: Collisions Mach number
15.6: Wave Behavior and Interaction kinetic theory of gases 16.3: Doppler Effect and Sonic Booms
26.2: Diffraction
26.3: Further Topics
12.1: Introduction machines
12.5: Kinetic Theory
26.4: Applications of Wave Optics 8.5: Applications of Statics
28.2: Applications of Quantum Mechanics magnetic domains
interferometer L 21.2: Magnets
26.1: Superposition and Interference laminar Magnetic Field
27.3: Relativistic Quantities 11.4: Other Applications 21.3: Magnetic Force on a Moving Electric Charge
intermolecular laminar flow 21.5: Magnetic Fields, Magnetic Forces, and
10.4: Cohesion and Adhesion 11.1: Overview Conductors
12.6: Phase Changes 21.6: Applications of Magnetism
land 22.1: Magnetic Flux, Induction, and Faraday’s Law
internal energy 26.4: Applications of Wave Optics 22.4: Magnetic Fields and Maxwell Revisited
12.4: Ideal Gas Law laser
13.1: Introduction
magnetic field lines
14.1: Introduction 28.2: Applications of Quantum Mechanics 21.1: Magnetism and Magnetic Fields
14.2: The First Law of Thermodynamics latent heat of fusion magnetic flux
14.3: The Second Law of Thermodynamics 13.3: Phase Change and Latent Heat 22.1: Magnetic Flux, Induction, and Faraday’s Law
International Systems of Units latent heat of vaporization magnetic mirror
12.4: Ideal Gas Law 13.3: Phase Change and Latent Heat 21.4: Motion of a Charged Particle in a Magnetic
inverse law Field
5.5: Newton’s Law of Universal Gravitation 1.1: The Basics of Physics magnetron
inviscid law of conservation of energy 21.4: Motion of a Charged Particle in a Magnetic
11.3: Bernoulli’s Equation 14.2: The First Law of Thermodynamics Field
ionization energy LCD Magnification
28.2: Applications of Quantum Mechanics 26.4: Applications of Wave Optics 24.3: Lenses
ionizing radiation length magnitude
23.1: The Electromagnetic Spectrum 1.2: Units 3.2: Vectors
30.4: Applications of Nuclear Physics
length contraction mass
iridescence 27.4: Implications of Special Relativity 21.6: Applications of Magnetism
26.2: Diffraction 4.2: Force and Mass
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mass distribution N optical window
15.3: Periodic Motion nanostructure 23.1: The Electromagnetic Spectrum
Mass Spectrometer 25.2: Other Optical Instruments order of magnitude
21.4: Motion of a Charged Particle in a Magnetic natural convection 1.3: Significant Figures and Order of Magnitude
Field origin
13.4: Methods of Heat Transfer
matter 3.2: Vectors
natural frequency
1.1: The Basics of Physics orthogonal
16.5: Further Topics
matter waves 17.4: The Electric Field Revisited
natural satellite
28.1: History and Quantum Mechanical Quantities 21.4: Motion of a Charged Particle in a Magnetic
29.2: The Early Atom 5.6: Kepler’s Laws Field
29.3: Atomic Physics and Quantum Mechanics nematic 26.1: Superposition and Interference
Maxwell’s Equations 26.4: Applications of Wave Optics oscillate
22.1: Magnetic Flux, Induction, and Faraday’s Law nerve impulse 15.2: Hooke’s Law
27.1: Introduction 16.4: Interactions with Sound Waves 15.3: Periodic Motion
28.1: History and Quantum Mechanical Quantities 15.7: Waves on Strings
29.2: The Early Atom
net force 26.3: Further Topics
4.3: Newton’s Laws
mean motion oscillator
5.6: Kepler’s Laws
neuron 15.4: Damped and Driven Oscillations
18.2: Equipotential Surfaces and Lines
mechanical advantage 18.3: Point Charge
osmosis
8.5: Applications of Statics 11.1: Overview
neurotransmitter
mechanical equivalent of heat 16.4: Interactions with Sound Waves
over damping
13.1: Introduction 15.4: Damped and Driven Oscillations
Newton’s third law of motion
media 7.4: Rocket Propulsion
oxidation
16.1: Introduction 14.2: The First Law of Thermodynamics
Newton's Law of Gravitation
medium 5.7: Gravitational Potential Energy
ozone
15.5: Waves 23.1: The Electromagnetic Spectrum
Newtonian mechanics
metabolism 12.5: Kinetic Theory
14.2: The First Law of Thermodynamics
NMR P
meteorology 30.4: Applications of Nuclear Physics parallel
5.4: Types of Forces in Nature 20.2: Resistors in Series and Parallel
noble gas
metric 12.5: Kinetic Theory parallel equivalent resistance
27.4: Implications of Special Relativity 19.3: Resistance and Resistors
node
microscopy 16.1: Introduction Paramagnetism
26.4: Applications of Wave Optics 16.5: Further Topics 21.6: Applications of Magnetism
microstate Nonconservative Forces particle accelerator
14.5: The Third Law of Thermodynamics 6.5: Potential Energy and Conservation of Energy 18.1: Overview
Minkowski space normal 28.2: Applications of Quantum Mechanics
27.4: Implications of Special Relativity 4.4: Other Examples of Forces pendentive
model normal force 8.5: Applications of Statics
1.1: The Basics of Physics 5.1: Introduction to UCM and Gravitation Pendulums
mole nuclear reactor 6.8: Further Topics
12.4: Ideal Gas Law 30.4: Applications of Nuclear Physics perihelion
moment of inertia nucleons 5.6: Kepler’s Laws
12.5: Kinetic Theory 30.3: Quantum Tunneling and Conservation Laws period
8.4: Solving Statics Problems 15.3: Periodic Motion
nucleosynthesis
Momentum 30.4: Applications of Nuclear Physics
16.1: Introduction
4.3: Newton’s Laws
nucleus periodic table
7.1: Introduction 29.5: Multielectron Atoms
7.3: Collisions 17.1: Overview
9.6: Conservation of Angular Momentum 29.2: The Early Atom permanent magnet
30.1: The Nucleus 21.1: Magnetism and Magnetic Fields
monochromatic
26.1: Superposition and Interference
nuclide Permeability
26.2: Diffraction 30.1: The Nucleus 22.1: Magnetic Flux, Induction, and Faraday’s Law
28.2: Applications of Quantum Mechanics null measurements permittivity
motion 20.4: Voltmeters and Ammeters 18.4: Capacitors and Dielectrics
2.4: Problem-Solving for Basic Kinematics perpendicular
Motional Emf O 4.4: Other Examples of Forces
4.7: Further Applications of Newton’s Laws
22.1: Magnetic Flux, Induction, and Faraday’s Law Objective measurement
MRI 16.1: Introduction
phase
23.2: Electromagnetic Waves and their Properties
30.4: Applications of Nuclear Physics ohm
muscles 19.2: Electric Current
phase velocity
15.5: Waves
8.5: Applications of Statics Ohm’s law
mutual inductance 22.2: AC Circuits
phasor
22.2: AC Circuits
22.2: AC Circuits ohmic
myocardium 19.3: Resistance and Resistors
Phasors
19.6: Electricity in the World 19.5: Alternating Currents
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phon Potential Energy rarefaction
16.2: Sound Intensity and Level 18.1: Overview 15.5: Waves
phosphor 5.7: Gravitational Potential Energy raster
5.8: Energy Conservation
18.5: Applications 18.5: Applications
6.8: Further Topics
photoconductivity potentiometer ray tracing
17.6: Applications of Electrostatics 24.3: Lenses
20.4: Voltmeters and Ammeters
Photoelectric effect power Rayleigh criterion
15.6: Wave Behavior and Interaction 28.1: History and Quantum Mechanical Quantities
15.5: Waves
28.1: History and Quantum Mechanical Quantities Reactance
6.6: Power
29.2: The Early Atom
photoelectron prefix 22.2: AC Circuits
1.2: Units reflection
28.1: History and Quantum Mechanical Quantities
photon pressure 24.1: Overview
24.2: Reflection, Refraction, and Dispersion
10.4: Cohesion and Adhesion
23.2: Electromagnetic Waves and their Properties 26.3: Further Topics
28.2: Applications of Quantum Mechanics Probability density function refraction
29.2: The Early Atom 28.1: History and Quantum Mechanical Quantities
24.1: Overview
physical pendulum propulsion 24.2: Reflection, Refraction, and Dispersion
15.3: Periodic Motion 5.8: Energy Conservation 25.2: Other Optical Instruments
pit pupil 26.3: Further Topics
26.4: Applications of Wave Optics 25.1: The Human Eye Refractive index
Planck constant Purkinje fibers 15.6: Wave Behavior and Interaction
23.2: Electromagnetic Waves and their Properties
29.1: Overview 19.6: Electricity in the World
Planck’s constant relative
3.4: Multiple Velocities
28.1: History and Quantum Mechanical Quantities Q
plane relativistic quantum theory
Quantization 28.1: History and Quantum Mechanical Quantities
10.6: Deformation of Solids 29.2: The Early Atom
Plane waves 29.3: Atomic Physics and Quantum Mechanics
relativity of simultaneity
27.4: Implications of Special Relativity
15.5: Waves quantum mechanics
26.1: Superposition and Interference 9.6: Conservation of Angular Momentum
renewable forms of energy
planet 6.7: CASE STUDY: World Energy Use
Quarks
5.1: Introduction to UCM and Gravitation 30.1: The Nucleus
reorientate
Plasma 3.3: Projectile Motion
10.1: Introduction R resisitivity
12.6: Phase Changes 17.1: Overview
radar 19.3: Resistance and Resistors
plumb line
23.1: The Electromagnetic Spectrum resistance
7.5: Center of Mass
radial 20.2: Resistors in Series and Parallel
point mass
18.1: Overview Resistor
5.5: Newton’s Law of Universal Gravitation 5.2: Non-Uniform Circular Motion
point particle 20.3: Kirchhoff’s Rules
Radians 20.5: RC Circuits
7.5: Center of Mass 5.3: Velocity, Acceleration, and Force
Poiseuille's Law resolution
radiation 26.2: Diffraction
10.2: Density and Pressure 1.2: Units
polarity resonance
radiative transfer 15.6: Wave Behavior and Interaction
17.6: Applications of Electrostatics 13.5: Global Warming 22.2: AC Circuits
polarizability radio waves rest mass
26.3: Further Topics 23.1: The Electromagnetic Spectrum 27.3: Relativistic Quantities
polarization 23.3: Applications of EM Waves restoring force
26.3: Further Topics Radioactive decay 15.3: Periodic Motion
polarizer 30.1: The Nucleus 15.6: Wave Behavior and Interaction
24.2: Reflection, Refraction, and Dispersion 30.2: Radioactivity resultant
30.4: Applications of Nuclear Physics
position 4.6: Vector Nature of Forces
radioactive tracer
2.3: Acceleration reversible
30.4: Applications of Nuclear Physics
positive feedback 14.2: The First Law of Thermodynamics
13.4: Methods of Heat Transfer
radiograph 14.4: Entropy
23.1: The Electromagnetic Spectrum Reynolds number
positron
27.3: Relativistic Quantities
radiography 11.2: Flow in Tubes
29.3: Atomic Physics and Quantum Mechanics 25.2: Other Optical Instruments right hand rule
30.2: Radioactivity radioisotope 21.3: Magnetic Force on a Moving Electric Charge
30.4: Applications of Nuclear Physics 30.2: Radioactivity 9.7: Vector Nature of Rotational Kinematics
potential radiometric dating rigid
12.3: Thermal Expansion 30.2: Radioactivity 4.7: Further Applications of Newton’s Laws
6.5: Potential Energy and Conservation of Energy
radionuclide rigid body
potential difference
30.1: The Nucleus 7.5: Center of Mass
18.1: Overview 30.2: Radioactivity
20.2: Resistors in Series and Parallel
radon
30.2: Radioactivity
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rms silver halide Static Equilibrium
12.5: Kinetic Theory 28.2: Applications of Quantum Mechanics 10.2: Density and Pressure
20.5: RC Circuits simple circuit 17.4: The Electric Field Revisited
22.2: AC Circuits 18.2: Equipotential Surfaces and Lines
19.3: Resistance and Resistors
rms current 18.3: Point Charge
Simple Harmonic Motion 8.3: Stability
19.5: Alternating Currents
rms voltage
15.6: Wave Behavior and Interaction stimulated emission
simple harmonic oscillator 29.4: Applications of Atomic Physics
19.5: Alternating Currents
root mean square
15.3: Periodic Motion Stoke's theorem
simple pendulum 22.1: Magnetic Flux, Induction, and Faraday’s Law
19.5: Alternating Currents
rotation
15.3: Periodic Motion strain
sinoatrial node 4.7: Further Applications of Newton’s Laws
9.6: Conservation of Angular Momentum
19.6: Electricity in the World 8.6: Elasticity, Stress, Strain, and Fracture
Rotational Inertia sinusoidal streamlined
8.8: Torque and Angular Acceleration
15.3: Periodic Motion 11.4: Other Applications
9.10: Conservation of Energy
9.4: Dynamics sinusoidal steady state stress
9.9: Linear and Rotational Quantities 19.5: Alternating Currents 12.8: Thermal Stresses
Snell’s law of refraction 4.7: Further Applications of Newton’s Laws
S 15.6: Wave Behavior and Interaction string theory
24.2: Reflection, Refraction, and Dispersion 28.1: History and Quantum Mechanical Quantities
scalar
2.1: Basics of Kinematics solenoid stroboscopic
3.2: Vectors 21.2: Magnets 2.4: Problem-Solving for Basic Kinematics
scanning tunneling microscope (STM) 22.1: Magnetic Flux, Induction, and Faraday’s Law Subjective measurement
28.2: Applications of Quantum Mechanics sonic boom 16.1: Introduction
Schrodinger's Equation 16.3: Doppler Effect and Sonic Booms sublimation
28.1: History and Quantum Mechanical Quantities Special relativity 10.1: Introduction
10.6: Deformation of Solids 13.3: Phase Change and Latent Heat
scientific method
1.1: The Basics of Physics
22.1: Magnetic Flux, Induction, and Faraday’s Law Superconductivity
23.2: Electromagnetic Waves and their Properties
19.3: Resistance and Resistors
scientific notation 26.4: Applications of Wave Optics
1.3: Significant Figures and Order of Magnitude 27.1: Introduction superimpose
27.2: Consequences of Special Relativity 16.4: Interactions with Sound Waves
Scintillation 27.3: Relativistic Quantities 25.2: Other Optical Instruments
29.2: The Early Atom 27.4: Implications of Special Relativity 26.4: Applications of Wave Optics
scintillators 28.1: History and Quantum Mechanical Quantities superposition
30.2: Radioactivity specific heat 15.6: Wave Behavior and Interaction
second law of thermodynamics 12.4: Ideal Gas Law 16.4: Interactions with Sound Waves
14.3: The Second Law of Thermodynamics 13.2: Specific Heat 17.4: The Electric Field Revisited
selective absorber spectral color surface vertices
13.5: Global Warming 23.1: The Electromagnetic Spectrum 24.3: Lenses
semiclassical spectral radiance symmetrical
29.2: The Early Atom 29.1: Overview 3.3: Projectile Motion
semiclassical approach spectrum symmetry
28.2: Applications of Quantum Mechanics 29.2: The Early Atom 4.3: Newton’s Laws
semiconductor speed of light systole
19.3: Resistance and Resistors 27.2: Consequences of Special Relativity 11.2: Flow in Tubes
30.2: Radioactivity 27.3: Relativistic Quantities
series speed of propagation T
15.5: Waves
20.2: Resistors in Series and Parallel tangent
series equivalent resistance spherical aberration 4.7: Further Applications of Newton’s Laws
25.2: Other Optical Instruments
19.3: Resistance and Resistors tangential acceleration
shear stress spin 9.1: Quantities of Rotational Kinematics
21.6: Applications of Magnetism
10.1: Introduction
28.2: Applications of Quantum Mechanics
telecommunication
11.2: Flow in Tubes
23.3: Applications of EM Waves
shock hazard stable equilibrium
8.3: Stability
temperature coefficient of resistivity
19.5: Alternating Currents
19.3: Resistance and Resistors
19.6: Electricity in the World standard atmosphere
Shunt resistance 12.2: Temperature and Temperature Scales
tensile stress
8.5: Applications of Statics
20.4: Voltmeters and Ammeters standing wave
shutter speed 15.7: Waves on Strings
terminal velocity
29.2: The Early Atom 17.1: Overview
25.2: Other Optical Instruments
SI units static terminal voltage
1.4: Solving Physics Problems 20.2: Resistors in Series and Parallel
12.4: Ideal Gas Law
sidereal year 4.7: Further Applications of Newton’s Laws tesla
5.6: Kepler’s Laws
static contrast ratio 21.3: Magnetic Force on a Moving Electric Charge
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theory Transformer vector field
1.1: The Basics of Physics 22.1: Magnetic Flux, Induction, and Faraday’s Law 10.6: Deformation of Solids
thermal agitation translation 17.4: The Electric Field Revisited
23.1: The Electromagnetic Spectrum 8.2: Conditions for Equilibrium velocity
thermal conductivity transmutation 2.2: Speed and Velocity
2.3: Acceleration
13.4: Methods of Heat Transfer 30.2: Radioactivity
3.2: Vectors
thermal energy transverse 4.2: Force and Mass
14.3: The Second Law of Thermodynamics 21.5: Magnetic Fields, Magnetic Forces, and 5.3: Velocity, Acceleration, and Force
thermal equilibrium Conductors ventricle
12.7: The Zeroth Law of Thermodynamics transverse wave 11.4: Other Applications
13.1: Introduction 15.5: Waves vessel
14.1: Introduction 15.7: Waves on Strings
11.2: Flow in Tubes
thermal hazard 16.5: Further Topics
trigonometry Vienna Standard Mean Ocean Water
19.5: Alternating Currents
12.4: Ideal Gas Law
19.6: Electricity in the World 1.4: Solving Physics Problems
thermal radiation triple point virtual image
24.4: Mirrors
23.1: The Electromagnetic Spectrum 12.2: Temperature and Temperature Scales
thermodynamic temperature trough viscosity
10.1: Introduction
14.1: Introduction 15.5: Waves
11.2: Flow in Tubes
Thermodynamics tunneling 11.3: Bernoulli’s Equation
12.2: Temperature and Temperature Scales 30.3: Quantum Tunneling and Conservation Laws 11.4: Other Applications
12.6: Phase Changes Turbine visible light
14.1: Introduction
22.1: Magnetic Flux, Induction, and Faraday’s Law 23.1: The Electromagnetic Spectrum
thermography turbulence voltage
23.1: The Electromagnetic Spectrum
11.4: Other Applications 17.1: Overview
Thick Lenses turbulent 19.2: Electric Current
24.3: Lenses
11.4: Other Applications
thin lens turbulent flow W
24.3: Lenses warfare
11.1: Overview
thin lens equation 30.4: Applications of Nuclear Physics
24.3: Lenses
U watt
Thomson scattering 6.6: Power
29.2: The Early Atom
Ultrasound
15.6: Wave Behavior and Interaction wave
Thoracic Cavity 15.5: Waves
10.2: Density and Pressure
underdamped
15.4: Damped and Driven Oscillations Wave equation
thrust 15.6: Wave Behavior and Interaction
4.3: Newton’s Laws
Uniform Circular Motion
15.3: Periodic Motion wave speed
time dilation 9.9: Linear and Rotational Quantities 15.5: Waves
27.2: Consequences of Special Relativity waveform
27.4: Implications of Special Relativity
uniform motion
4.3: Newton’s Laws 15.1: Introduction
tomography
unit vector wavefront
25.2: Other Optical Instruments
30.4: Applications of Nuclear Physics 17.3: Coulomb’s Law 16.5: Further Topics
torque 3.2: Vectors wavelength
10.2: Density and Pressure 15.5: Waves
18.1: Overview V 17.2: Shelding and Charging Through Induction
22.3: Applications of Induction and EM Waves
21.5: Magnetic Fields, Magnetic Forces, and valance shell 23.2: Electromagnetic Waves and their Properties
Conductors
29.5: Multielectron Atoms 26.1: Superposition and Interference
22.1: Magnetic Flux, Induction, and Faraday’s Law
6.4: Work-Energy Theorem vapor 26.3: Further Topics
7.5: Center of Mass 12.6: Phase Changes weight
8.2: Conditions for Equilibrium vaporization 5.5: Newton’s Law of Universal Gravitation
8.4: Solving Statics Problems
8.8: Torque and Angular Acceleration
13.6: Phase Equilbrium wettability
9.4: Dynamics vascular 11.3: Bernoulli’s Equation
9.5: Rotational Kinetic Energy 11.1: Overview Wheatstone bridge
9.6: Conservation of Angular Momentum vector 20.4: Voltmeters and Ammeters
9.7: Vector Nature of Rotational Kinematics
9.9: Linear and Rotational Quantities 17.4: The Electric Field Revisited work
18.1: Overview 15.5: Waves
torques 2.1: Basics of Kinematics 18.1: Overview
4.7: Further Applications of Newton’s Laws 3.2: Vectors 6.2: Work Done by a Constant Force
Torr 4.2: Force and Mass 6.3: Work Done by a Variable Force
4.6: Vector Nature of Forces 9.10: Conservation of Energy
10.2: Density and Pressure
5.9: Angular vs. Linear Quantities
tracer 8.1: Introduction
30.4: Applications of Nuclear Physics vector area X
trajectory 22.1: Magnetic Flux, Induction, and Faraday’s Law xerography
3.3: Projectile Motion 17.6: Applications of Electrostatics
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Z
zeroth law of thermodynamics
12.7: The Zeroth Law of Thermodynamics
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CHAPTER OVERVIEW
This page titled 12: Temperature and Kinetic Theory is shared under a not declared license and was authored, remixed, and/or curated by
Boundless.
1
12.1: Introduction
learning objectives
Describe gas using the kinetic theory of gases
Translational Motion of Helium: Real gases do not always behave according to the ideal model under certain conditions, such as
high pressure. Here, the size of helium atoms relative to their spacing is shown to scale under 1950 atmospheres of pressure.
The kinetic theory of gases uses the model of the ideal gas to relate temperature to the average translational kinetic energy of the
molecules in a container of gas in thermodynamic equilibrium. Classical mechanics defines the translational kinetic energy of a gas
molecule as follows:
1
2
Ek = mv , (12.1.1)
2
where m is the particle mass and v its speed (the magnitude of its velocity). The distribution of the speeds (which determine the
translational kinetic energies) of the particles in a classical ideal gas is called the Maxwell-Boltzmann distribution. In kinetic
theory, the temperature of a classical ideal gas is related to its average kinetic energy per degree of freedom Ek via the equation:
1
¯
Ek = kT, (12.1.2)
2
(k: Boltzmann’s constant). We will derive this relationship in the following atoms. We will also derive the ideal gas law:
pV = nRT, (12.1.3)
(R: ideal gas constant, n: number of moles of gas) from a microscopic theory.
learning objectives
Formulate five postulates of John Dalton’s atomic theory
Atomic theory is a scientific theory of the nature of matter which states that matter is composed of discrete units called atoms, as
opposed to the obsolete notion that matter could be divided into any arbitrarily small quantity. Although physicists discovered that
the so-called “indivisible atom” was actually a conglomerate of various subatomic particles, the concept of atoms is still important
because they are building blocks of matter and form the basis of chemistry.
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Illustration of the Helium Atom: This is an illustration of the helium atom, depicting the nucleus (pink) and the electron cloud
distribution (black). The nucleus (upper right) in helium-4 is in reality spherically symmetric and closely resembles the electron
cloud, although for more complicated nuclei this is not always the case. The black bar is one angstrom (10-10 m, or 100 pm).
5 Main Points
Dalton’s atomic theory had 5 main points:
1. Elements are made of extremely small particles called atoms.
2. Atoms of a given element are identical in size, mass, and other properties; atoms of different elements differ in size, mass, and
other properties.
3. Atoms cannot be subdivided, created, or destroyed.
4. Atoms of different elements combine in simple whole-number ratios to form chemical compounds.
5. In chemical reactions, atoms are combined, separated, or rearranged.
Of these five, only three are still considered valid today. 1, 4, and 5 are valid, while 2 and 3 have turned out not to be the case.
Atoms can be broken down into smaller pieces, and atoms of a given element can vary in mass and other properties (see isotopes
and ions).
Knowing that a gas is composed of small atomic and molecular particles, it is natural to try to explain properties of the gas from a
microscopic point of view. This effort led to the development of the kinetic theory of gases, where macroscopic properties of gases,
such as pressure, temperature, and volume, are explained by considering their molecular composition and motion.
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John Dalton’s A New System of Chemical Philosophy: Various atoms and molecules as depicted in John Dalton’s A New System
of Chemical Philosophy (1808).
Key Points
The kinetic theory posits that pressure is due to collisions between molecules moving at different velocities through Brownian
motion.
The temperature of an ideal monatomic gas is a measure of the average kinetic energy of its atoms. In kinetic theory, it is related
1
to its average kinetic energy per degree of freedom Ek via the equation: E
¯
k = kT .
2
The kinetic theory of gases uses the model of the ideal gas to relate temperature to the average translational kinetic energy of
the molecules in a container of gas in thermodynamic equilibrium.
John Dalton was the first to propose a scientific theory of atoms. He based his study on two laws: the law of conservation of
mass and the law of definite proportions.
Dalton proposed that each chemical element is composed of atoms of a single, unique type, and though they cannot be altered
or destroyed by chemical means, they can combine to form more complex structures.
Kinetic theory of gases explain macroscopic properties of gases, such as pressure, temperature, and volume, by considering
their molecular composition and motion.
While Dalton’s idea of matter being composed of various atoms was correct, he was wrong about some of their properties.
Atoms can be broken down into smaller parts. Atoms of the same element can have slightly different masses and behave
differently. See isotopes and ions for examples.
Key Terms
ideal gas: A hypothetical gas whose molecules exhibit no interaction and undergo elastic collision with each other and with the
walls of the container.
degree of freedom: Any of the coordinates, a minimum number of which are needed to specify the motion of a mechanical
system.
Brownian motion: Random motion of particles suspended in a fluid, arising from those particles being struck by individual
molecules of the fluid.
atom: The smallest possible amount of matter which still retains its identity as a chemical element, now known to consist of a
nucleus surrounded by electrons.
kinetic theory of gases: The kinetic theory of gases describes a gas as a large number of small particles (atoms or molecules),
all of which are in constant, random motion.
chemical reaction: A process, involving the breaking or making of interatomic bonds, in which one or more substances are
changed into others.
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LICENSES AND ATTRIBUTIONS
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12.2: Temperature and Temperature Scales
learning objectives
Explain how the Celsius scale is defined
Celsius, also known as centigrade, is a scale to measure temperature. The unit of measurement is the degree Celsius (°C). It is one
of the most commonly used temperature units in the world. The unit system is named after the Swedish astronomer Anders Celsius
(1701-1744), who developed a similar temperature scale.
AP Physics 2 - Temperature
Temperature Scales: A brief introduction to temperature and temperature scales for students studying thermal physics or
thermodynamics.
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Thermometer: A thermometer calibrated in degrees Celsius
From 1743 until 1954, 0°C was defined as the freezing point of water, and 100°C was defined as the boiling point of water, both at
a pressure of one standard atmosphere, with mercury as the working material. Although these defining correlations are commonly
taught in schools today, by international agreement the unit “degree Celsius” and the Celsius scale are currently defined by two
different temperatures: absolute zero and the triple point of Vienna Standard Mean Ocean Water (VSMOW; specially purified
water). This definition also precisely relates the Celsius scale to the Kelvin scale, which defines the SI base unit of thermodynamic
temperature and which uses the symbol K. Absolute zero, the lowest temperature possible (the temperature at which matter reaches
minimum entropy), is defined as being precisely 0K and -273.15°C. The temperature of the triple point of water is defined as
precisely 273.16K and 0.01°C. Based on this, the relationship between degree Celsius and Kelvin is as follows:
Phase Diagram of Water: In this typical phase diagram of water, the green lines mark the freezing point, and the blue line marks
the boiling point, showing how they vary with pressure. The dotted line illustrates the anomalous behavior of water. Note that water
changes states based on the pressure and temperature.
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TCelsius = TKelvin − 273.15. (12.2.1)
Besides expressing specific temperatures along its scale (e.g., “Gallium melts at 29.7646°C” and “The temperature outside is 23
degrees Celsius”), the degree Celsius is also suitable for expressing temperature intervals — differences between temperatures, or
their uncertainties (e.g. “The output of the heat exchanger is hotter by 40 degrees Celsius” and “Our standard uncertainty is ±3°C”).
Because of this dual usage, one must not rely upon the unit name or its symbol to denote that a quantity is a temperature interval; it
must be clear through context or explicit statement that the quantity is an interval.
Farenheit Scale
In the Fahrenheit scale, the freezing of water is defined at 32 degrees, while the boiling point of water is defined to be 212 degrees.
learning objectives
Explain how the Fahrenheit scale is defined and convert between it and Celsius
The Fahrenheit scale measures temperature. It is based on a scale proposed in 1724 by physicist Daniel Gabriel Fahrenheit (1686-
1736). The unit of this scale is the degree Fahrenheit (°F). On this scale, water’s freezing point is defined to be 32 degrees, while
water’s boiling point is defined to be 212 degrees.
Historically, the zero point of the Fahrenheit scale was determined by evaluating a thermometer placed in brine. Fahrenheit himself
used a mixture of ice, water, and ammonium chloride (a salt) at a 1:1:1 ratio. This is a frigorific mixture, which stabilizes its
temperature automatically; the stable temperature of this mixture was defined as 0 °F (-17.78 °C). The second determining point,
32 degrees, was a mixture of just ice and water at a 1:1 ratio. The third determining point, 96 degrees, was approximately the
temperature of the human body, then called “blood-heat. ”
The Fahrenheit system puts the boiling and freezing points of water exactly 180 degrees apart. Therefore, a degree on the
Fahrenheit scale is 1/180 of the interval between the freezing point and the boiling point. On the Celsius scale, the freezing and
boiling points of water are 100 degrees apart. A temperature interval of 1 °F is equal to an interval of 5/9 degrees Celsius (°C). To
convert °F to °C, you can use the following formula:
5
TCelsius = (TFahrenheit − 32) (12.2.2)
9
The Fahrenheit and Celsius scales intersect at -40° (-40 °F and -40 °C represent the same temperature). Absolute zero (-273.15 °C,
or 0K) is defined as -459.67 °F.
The Fahrenheit scale was replaced by the Celsius scale in most countries in the mid- to late-20th century, though Canada retains it
as a supplementary scale that can be used alongside the Celsius scale. The Fahrenheit scale remains the official scale of the United
States, the Cayman Islands, Palau, the Bahamas, and Belize.
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Fig 2: Comparison of Celsius vs Fahrenheit scales.
AP Physics 2 - Temperature
12.2.4 https://phys.libretexts.org/@go/page/14509
Temperature Scales: A brief introduction to temperature and temperature scales for students studying thermal physics or
thermodynamics.
Absolute Zero
Absolute zero is the coldest possible temperature; formally, it is the temperature at which entropy reaches its minimum value.
learning objectives
Explain why absolute zero is a natural choice as the null point for a temperature unit system
Absolute zerois the coldest possible temperature. Formally, it is the temperature at which entropy reaches its minimum value. More
simply put, absolute zero refers to a state in which all the energy of a system is extracted (by definition, the lowest energy state the
system can have). Absolute zero is universal in the sense that all matteris in ground state at this temperature. Therefore, it is a
natural choice as the null point for a temperature unit system.
Graph of Pressure Versus Temperature: Graph of pressure versus temperature for various gases kept at a constant volume. Note
that all of the graphs extrapolate to zero pressure at the same temperature
To be precise, a system at absolute zero still possesses quantum mechanical zero-point energy, the energy of its ground state. The
uncertainty principle states that the position of a particle cannot be determined with absolute precision; therefore a particle is in
motion even if it is at absolute zero, and a ground state still carries a minimal amount of kinetic energy. However, in the
interpretation of classical thermodynamics, kinetic energy can be zero, and the thermal energy of matter vanishes.
The zero point of a thermodynamic temperature scale, such as the Kelvin scale, is set at absolute zero. By international agreement,
absolute zero is defined as 0K on the Kelvin scale and as -273.15° on the Celsius scale (equivalent to -459.67° on the Fahrenheit
scale). Scientists have brought systems to temperatures very close to absolute zero, at which point matter exhibits quantum effects
such as superconductivity and superfluidity. The lowest temperature that has been achieved in the laboratory is in the 100 pK range,
where pK (pico-Kelvin) is equivalent to 10-12 K. The lowest natural temperature ever recorded is approximately 1K, observed in
the rapid expansion of gases leaving the Boomerang Nebula, shown below.
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Boomerang Nebula: The rapid expansion of gases resulting in the Boomerang Nebula causes the lowest observed temperature
outside a laboratory.
AP Physics 2 - Temperature
Temperature Scales: A brief introduction to temperature and temperature scales for students studying thermal physics or
thermodynamics.
Kelvin Scale
The kelvin is a unit of measurement for temperature; the null point of the Kelvin scale is absolute zero, the lowest possible
temperature.
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learning objectives
Explain how the Kelvin scale is defined
The kelvin is a unit of measurement for temperature. It is one of the seven base units in the International System of Units (SI) and
is assigned the unit symbol K. The Kelvin scale is an absolute, thermodynamic temperature scale using absolute zero as its null
point. In the classical description of thermodynamics, absolute zero is the temperature at which all thermal motion ceases.
The choice of absolute zero as null point for the Kelvin scale is logical. Different types of matter boil or freeze at different
temperatures, but at 0K (absolute zero), all thermal motions of any matter are maximally suppressed. The Kelvin scale is used
extensively in scientific work because a number of physical quantities, such as the volume of an ideal gas, are directly related to
absolute temperature.
The Kelvin scale is named after Glasgow University engineer and physicist William Thomson, 1st Baron Kelvin (1824-1907), who
wrote of the need for an “absolute thermometric scale. ” Unlike the degree Fahrenheit and the degree Celsius, the kelvin is not
referred to or typeset as a degree. The kelvin is the primary unit of measurement in the physical sciences, but it is often used in
conjunction with the degree Celsius, which has the same magnitude. The kelvin is defined as the fraction 1/273.16 of the
thermodynamic temperature of the triple point of water (exactly 0.01°C, or 32.018°F). To convert kelvin to degrees Celsius, we use
the following formula:
Subtracting 273.16K from the temperature of the triple point of water, 0.01°C, makes absolute zero (0K) equivalent to -273.15°C
and -460°F.
Calculating U
Relationships Between the Temperature Scales: Relationships between the Fahrenheit, Celsius, and Kelvin temperature scales,
rounded to the nearest degree. The relative sizes of the scales are also shown
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AP Physics 2 - Temperature
Temperature Scales: A brief introduction to temperature and temperature scales for students studying thermal physics or
thermodynamics.
Key Points
The degree Celsius (°C) can refer to a specific temperature on the Celsius scale as well as a unit to indicate a temperature
interval, a difference between two temperatures or an uncertainty.
The Celsius scale is currently defined by two different temperatures: absolute zero and the triple point of Vienna Standard Mean
Ocean Water (VSMOW; specially purified water).
Based on this, the relationship between Celsius and Kelvin is as follows: T =T
Celsius
− 273.15 .
Kelvin
The Fahrenheit system puts the boiling and freezing points of water exactly 180 degrees apart. Therefore, a degree on the
Fahrenheit scale is 1⁄180 of the interval between the freezing point and the boiling point.
5
To convert °F to °C, you can use the following formula: T Celsius = (TFahrenheit − 32) . The Fahrenheit and Celsius scales
9
intersect at −40°.
The Fahrenheit scale was replaced by the Celsius scale in most countries during the mid to late 20th century. Fahrenheit remains
the official scale of the United States, Cayman Islands, Palau, Bahamas and Belize.
Absolute zero is universal in the sense that all matter is in ground state at this temperature. Therefore, it is a natural choice as
the null point for a temperature unit system.
K system at absolute zero still possesses quantum mechanical zero-point energy, the energy of its ground state. However, in the
interpretation of classical thermodynamics, kinetic energy can be zero, and the thermal energy of matter vanishes.
The lowest temperature that has been achieved in the laboratory is in the 100 pK range, where pK (pico- Kelvin ) is equivalent
to 10-12 K. The lowest natural temperature ever recorded is approximately 1K, observed in the rapid expansion of gases leaving
12.2.8 https://phys.libretexts.org/@go/page/14509
the Boomerang Nebula.
0K ( absolute zero ) is universal because all thermal motions of all matter are maximally suppressed at this temperature.
Absolute zero is therefore the natural choice as the null point of the Kelvin scale.
The Kelvin scale is used extensively in scientific work because a number of physical quantities, such as the volume of an ideal
gas, are directly related to absolute temperature.
To convert kelvin to degree Celsius, we use the following formula: T Celsius =TKelvin− 273.15 .
Key Terms
kelvin: in the International System of Units, the base unit of thermodynamic temperature; 1/273.16 of the thermodynamic
temperature of the triple point of water; symbolized as K
absolute zero: The coldest possible temperature: zero on the Kelvin scale and approximately -273.15°C and -459.67°F. The
total absence of heat; the temperature at which motion of all molecules would cease.
standard atmosphere: an international reference pressure defined as 101.325 kPa and formerly used as a unit of pressure
brine: a solution of salt (usually sodium chloride) in water
frigorific mixture: A mixture of two or more chemicals that reaches an equilibrium temperature independent of the
temperature of any of its constituent chemicals. The temperature is also relatively independent of the quantities of mixtures as
long as a significant amount of each original chemical is present in its pure form
entropy: A measure of how evenly energy (or some analogous property) is distributed in a system.
thermodynamics: a branch of natural science concerned with heat and its relation to energy and work
absolute zero: The coldest possible temperature: zero on the Kelvin scale and approximately -273.15°C and -459.67°F. The
total absence of heat; the temperature at which motion of all molecules would cease.
Triple point: The unique temperature and pressure at which the solid, liquid and gas phases of a substance are all in
equilibrium.
ideal gas: A hypothetical gas whose molecules exhibit no interaction and undergo elastic collision with each other and with the
walls of the container.
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12.3: Thermal Expansion
learning objectives
Describe volume changes that take place in response to a temperature change
Thermal expansion is the tendency of matter to change in volume in response to a change in temperature. (An example of this is the
buckling of railroad track, as seen in. ) Atoms and molecules in a solid, for instance, constantly oscillate around its equilibrium
point. This kind of excitation is called thermal motion. When a substance is heated, its constituent particles begin moving more,
thus maintaining a greater average separation with their neighboring particles. The degree of expansion divided by the change in
temperature is called the material’s coefficient of thermal expansion; it generally varies with temperature.
Fig 1: Thermal expansion of long continuous sections of rail tracks is the driving force for rail buckling. This phenomenon resulted
in 190 train derailments during 1998–2002 in the US alone.
12.3.1 https://phys.libretexts.org/@go/page/14510
Thermal Expansion: A brief introduction to thermal expansion for students.
Linear Expansion
To a first approximation, the change in length measurements of an object (lineardimension as opposed to, for example, volumetric
dimension) due to thermal expansion is related to temperature change by a linear expansion coefficient. It is the fractional change
in length per degree of temperature change. Assuming negligible effect of pressure, we may write:
1 dL
αL = , (12.3.1)
L dT
where L is a particular length measurement and is the rate of change of that linear dimension per unit change in temperature.
dL
dT
From the definition of the expansion coefficient, the change in the linear dimension ΔL over a temperature range ΔT can be
estimated to be:
ΔL
= αL ΔT. (12.3.2)
L
This equation works well as long as the linear-expansion coefficient does not change much over the change in temperature. If it
does, the equation must be integrated.
Area Expansion
Objects expand in all dimensions. That is, their areas and volumes, as well as their lengths, increase with temperature.
learning objectives
Express the area thermal expansion coefficient in the form of an equation
We learned about the linear expansion (in one dimension) in the previous Atom. Objects expand in all dimensions, and we can
extend the thermal expansion for 1D to two (or three) dimensions. That is, their areas and volumes, as well as their lengths,
increase with temperature.
Quiz
Before we look into details, here is an interesting question. Imagine that we have a rectangular sheet of metal with a circular hole in
the middle. If the metal is heated, we can guess that the the piece, in general, will get larger due to thermal expansion. Now, what is
going to happen with the circular hole in the middle? Is the hole going to be larger or smaller? Answer: Imagine that we have a
12.3.2 https://phys.libretexts.org/@go/page/14510
similar metal sheet but without a hole. Draw an imaginary circular line representing the circular hole in our quiz. How does this
imaginary circle change as the metal is heated? Yes. It will get bigger. Therefore, you can guess that the hole in our quiz will get
larger.
Fig 1: In general, objects expand in all directions as temperature increases. In these drawings, the original boundaries of the objects
are shown with solid lines, and the expanded boundaries with dashed lines. (a) Area increases because both length and width
increase. The area of a circular plug also increases. (b) If the plug is removed, the hole it leaves becomes larger with increasing
temperature, just as if the expanding plug were still in place.
dT
is the rate of change of that area per unit change in temperature. The change in the linear dimension
can be estimated as: ΔA
A
= α ΔT . This equation works well as long as the linear expansion coefficient does not change much
A
over the change in temperature ΔT. If it does, the equation must be integrated.
ΔL
= A + 2A (12.3.6)
L
ΔA ΔL
The approximation holds for a sufficiently small ΔL campared to L. Since =2 from the equation above (and from the
A L
definitions of the thermal coefficients), we get α
A
= 2αL .
Volume Expansion
Substances expand or contract when their temperature changes, with expansion or contraction occurring in all directions.
learning objectives
Compare the effects of the pressure on the expansion of gaseous and solid materials
The volumetric thermal expansion coefficient is the most basic thermal expansion coefficient. illustrates that, in general, substances
expand or contract when their temperature changes, with expansion or contraction occurring in all directions. Such substances that
expand in all directions are called isotropic. For isotropic materials, the area and linear coefficients may be calculated from the
volumetric coefficient (discussed below).
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Volumetric Expansion: In general, objects expand in all directions as temperature increases. In these drawings, the original
boundaries of the objects are shown with solid lines, and the expanded boundaries with dashed lines. (a) Area increases because
both length and width increase. The area of a circular plug also increases. (b) If the plug is removed, the hole it leaves becomes
larger with increasing temperature, just as if the expanding plug were still in place. (c) Volume also increases, because all three
dimensions increase.
Thermal Expansion – Volume Expansion: A brief introduction to thermal expansion for students.
Mathematical definitions of these coefficients are defined below for solids, liquids, and gasses:
1 ∂V
αV = ( )p. (12.3.7)
V ∂T
The subscript p indicates that the pressure is held constant during the expansion. In the case of a gas, the fact that the pressure is
held constant is important, as the volume of a gas will vary appreciably with pressure as well as with temperature.
For a solid, we can ignore the effects of pressure on the material, thus the volumetric thermal expansion coefficient can be written:
12.3.4 https://phys.libretexts.org/@go/page/14510
1 dV
αV = , (12.3.8)
V dT
where V is the volume of the material, and is dV/dT the rate of change of that volume with temperature. This means that the
volume of a material changes by some fixed fractional amount. For example, a steel block with a volume of 1 cubic meter might
expand to 1.002 cubic meters when the temperature is raised by 50 °C. This is an expansion of 0.2%. The volumetric expansion
coefficient would be 0.2% for 50 °C, or 0.004% per degree C.
3 2 2 3
=L + 3 L ΔL + 3L(ΔL) + (ΔL) (12.3.10)
3 2
≈L + 3 L ΔL (12.3.11)
ΔL
= V + 3V . (12.3.12)
L
learning objectives
Describe thermal expansion properties of water
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Water Density vs. Temperature: The density of water as a function of temperature. Note that the thermal expansion is actually
very small. The maximum density at +4ºC is only 0.0075% greater than the density at 2ºC, and 0.012% greater than that at 0ºC.
Perhaps the most striking effect of this phenomenon is the freezing of water in a pond. When water near the surface cools down to
4ºC it is denser than the remaining water and thus will sink to the bottom. This “turnover” results in a layer of warmer water near
the surface, which is then cooled. Eventually the pond has a uniform temperature of 4ºC. If the temperature in the surface layer
drops below 4ºC, the water is less dense than the water below, and thus stays near the top.
As a result, the pond surface can completely freeze over, while the bottom may remain at 4ºC. The ice on top of liquid water
provides an insulating layer from winter’s harsh exterior air temperatures. Fish and other aquatic life can survive in 4ºC water
beneath ice, due to this unusual characteristic of water. It also produces circulation of water in the pond that is necessary for a
healthy ecosystem of the body of water.
Temperature in a Lake: Temperature distribution in a lake on warm and cold days in winter
Key Points
Inter-particle potential usually takes an asymmetric form, rather than a symmetric form as a function of particle-particle
distance. This is why matters expands and contracts as temperature changes.
The change in length measurements of an object due to thermal expansion is related to temperature change by a “linear
1 dL
expansion coefficient”, which is given as α L = .
L dT
The linear expansion coefficient is as an approximation over a narrow temperature interval only.
The area thermal expansion coefficient relates the change in a material’s area dimensions to a change in temperature. It is
1 dA
defined as α A = .
A dT
The relationship between the area and linear thermal expansion coefficient is given as the following: α = 2α .
A L
Just like the linear expansion coefficient, the area thermal expansion coefficient works as an approximation over a narrow
temperature interval only.
Substances that expand at the same rate in every direction are called isotropic.
In the case of a gas, expansion depends on how the pressure changed in the process because the volume of a gas will vary
appreciably with pressure as well as temperature.
For a solid, we can ignore the effects of pressure on the material, and the volumetric thermal expansion coefficient can be
1 dV
written as αV = . For isotropic materials, αV = 3 αL .
V dT
Water expands with increasing temperature (its density decreases) when it is at temperatures greater than 4ºC (40ºF). However,
it expands with decreasing temperature when it is between +4ºC and 0ºC (40ºF to 32ºF). Water is densest at +4ºC.
Due to the peculiar thermal expansion property of water, a pond surface can completely freeze over, while the bottom may
remain at 4ºC. Fish and other aquatic life can survive in 4ºC water beneath ice, due to this unusual characteristic of water.
The solid form of most substances is denser than the liquid phase; thus, a block of most solids will sink in the liquid. However,
a block of ice floats in liquid water because ice is less dense.
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Key Terms
potential: A curve describing the situation where the difference in the potential energies of an object in two different positions
depends only on those positions.
linear thermal expansion coefficient: The fractional change in length per degree of temperature change.
isotropic: Having properties that are identical in all directions; exhibiting isotropy.
hydrogen bond: A weak bond in which a hydrogen atom in one molecule is attracted to an electronegative atom (usually
nitrogen or oxygen) in the same or different molecule.
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12.4: Ideal Gas Law
learning objectives
Describe how ideal gas law was derived.
The ideal gas law is the equation of state of a hypothetical ideal gas (an illustration is offered in ). In an ideal gas, there is no
molecule-molecule interaction, and only elastic collisions are allowed. It is a good approximation to the behavior of many gases
under many conditions, although it has several limitations. It was first stated by Émile Clapeyron in 1834 as a combination of
Boyle’s law and Charles’ law.
Atoms and Modules in a Gas: Atoms and molecules in a gas are typically widely separated, as shown. Because the forces
between them are quite weak at these distances, they are often described by the ideal gas law.
Empirical Derivation
Boyle’s law states that pressure P and volume V of a given mass of confined gas are inversely proportional:
1
P ∝ , (12.4.1)
V
while Charles’ law states that volume of a gas is proportional to the absolute temperature T of the gas at constant pressure
V ∝ T. (12.4.2)
where C is a constant which is directly proportional to the amount of gas, n (representing the number of moles).
The proportionality factor is the universal gas constant, R, i.e. C = nR .
Hence the ideal gas law
PV = nRT (12.4.4)
Microscopic version
We have seen in the Atom on “Origin of Pressure” that
¯
Nmv2
P = , (12.4.5)
3V
where P is the pressure, N is the number of molecules, m is the mass of the molecule, v is the speed of molecules, and V is the
volume of the gas. Therefore, we derive a microscopic version of the ideal gas law
12.4.1 https://phys.libretexts.org/@go/page/14511
1 ¯2
PV = Nmv (12.4.6)
3
Isotherms
An isothermal process is a change of a system in which the temperature remains constant: ΔT = 0 .
learning objectives
Identify conditions at which isothermal processes can occur.
An isothermal process is a change of a system in which the temperature remains constant: ΔT = 0. Typically this occurs when a
system is in contact with an outside thermal reservoir (heat bath), and the change occurs slowly enough to allow the system to
adjust continually to the temperature of the reservoir through heat exchange. In contrast, an adiabatic process occurs when a system
exchanges no heat with its surroundings (Q = 0). In other words, in an isothermal process, the value ΔT = 0 but Q ≠ 0, while in an
adiabatic process, ΔT ≠ 0 but Q = 0.
For an ideal gas, the product PV (P: pressure, V: volume) is a constant if the gas is kept at isothermal conditions (Boyle’s law).
According to the ideal gas law, the value of the constant is NkT, where N is the number of molecules of gas and k is Boltzmann’s
constant.
This means that p = NkT
V
=
Constant
V
holds.
The family of curves generated by this equation is shown in the graph presented in. Each curve is called an isotherm. Such graphs
are termed indicator diagrams—first used by James Watt and others to monitor the efficiency of engines. The temperature
corresponding to each curve in the figure increases from the lower left to the upper right.
Calculation of Work
In thermodynamics, the work involved when a gas changes from state A to state B is simply:
VB
(This equation is derived in our Atom on “Constant Pressure” under kinetic theory. Note that P = F
A
. This definition is consistent
with our definition of work being force times distance. )
For an isothermal, reversible process, this integral equals the area under the relevant pressure-volume isotherm, and is indicated in
blue in for an ideal gas. Again, P = applies and with T being constant (as this is an isothermal process), we have:
nRT
12.4.2 https://phys.libretexts.org/@go/page/14511
Work Done by Gas During Expansion: The blue area represents “work” done by the gas during expansion for this isothermal
change.
VB VB
NkT
WA→B = ∫ pdV = ∫ dV (12.4.8)
VA VA
V
VB
= NkT ln . (12.4.9)
VA
By convention, work is defined as the work the system does on its environment. If, for example, the system expands by a piston
moving in the direction of force applied by the internal pressure of a gas, then the work is counted as positive. As this work is done
by using internal energy of the system, the result is that the internal energy decreases. Conversely, if the environment does work on
the system so that its internal energy increases, the work is counted as negative (for details on internal energy, check our Atom on
“Internal Energy of an Ideal Gas”).
Constant Pressure
Isobaric processis a thermodynamic process in which the pressure stays constant (at constant pressure, work done by a gas is
PΔV).
learning objectives
Describe behavior of monatomic gas during isobaric processes.
Under a certain constraint (e.g., pressure), gases can expand or contract; depending on the type of constraint, the final state of the
gas may change. For example, an ideal gas that expands while its temperature is kept constant (called isothermal process) will exist
in a different state than a gas that expands while pressure stays constant (called isobaric process). This Atom addresses isobaric
process and correlated terms. We will discuss isothermal process in a subsequent Atom.
Isobaric Process
An isobaric process is a thermodynamic process in which pressure stays constant: ΔP = 0 . For an ideal gas, this means the volume
of a gas is proportional to its temperature (historically, this is called Charles’ law ). Let’s consider a case in which a gas does work
on a piston at constant pressure P, referring to Fig 1 as illustration. Since the pressure is constant, the force exerted is constant and
the work done is given as W=Fd, where F (=PA) is the force on the piston applied by the pressure and d is the displacement of the
piston. Therefore, the work done by the gas (W) is:
W = PAd (12.4.10)
Because the change in volume of a cylinder is its cross-sectional area A times the displacement d, we see that Ad=ΔV, the change
in volume. Thus,
W = PΔV (12.4.11)
(as seen in Fig 2—isobaric process ). Note: if ΔV is positive, then W is positive, meaning that work is done by the gas on the
outside world. Using the ideal gas law PV = NkT(P = const) ,
12.4.3 https://phys.libretexts.org/@go/page/14511
Fig 2: A graph of pressure versus volume for a constant-pressure, or isobaric process. The area under the curve equals the work
done by the gas, since W=PΔV.
W = NkΔT (12.4.12)
Monatomic Gas
According to the first law of thermodynamics,
Q = ΔU + W (12.4.13)
(Eq. 2), where W is work done by the system, U is internal energy, and Q is heat. The law says that the heat transferred to the
system does work but also changes the internal energy of the system. Since,
U =
3
2
NkT for a monatomic gas, we get ΔU = 3
2
NkΔT
(Eq. 3; for the details on internal energy, see our Atom on “Internal Energy of an Ideal Gas”). By using the Equations 1 and 3, Eq. 2
can be written as:
Q=
5
2
NkΔT for monatomic gas in an isobaric process.
Specific Heat
Specific heat at constant pressure is defined by the following equation:
Q = ncPΔT
Here n is the amount of particles in a gas represented in moles. By noting that N = N An and R = kN (NA: Avogadro’s number,
A
2
kNA =
5
2
R for a monatomic gas.
Problem Solving
With the ideal gas law we can figure pressure, volume or temperature, and the number of moles of gases under ideal
thermodynamic conditions.
learning objectives
Identify steps used to solve the ideal gas equation.
The Ideal Gas Law is the equation of state of a hypothetical ideal gas. It is a good approximation to the behavior of many gases
under many conditions, although it has several limitations. It is most accurate for monatomic gases at high temperatures and low
pressures.
The ideal gas law has the form:
PV = nRT, (12.4.14)
12.4.4 https://phys.libretexts.org/@go/page/14511
where R is the universal gas constant, and with it we can find values of the pressure P, volume V, temperature T, or number of
moles n under a certain ideal thermodynamic condition. Typically, you are given enough parameters to calculate the unknown.
Variations of the ideal gas equation may help solving the problem easily. Here are some general tips.
The ideal gas law can also come in the form:
PV = NkT, (12.4.15)
where N is the number of particles in the gas and k is the Boltzmann constant.
To solve the ideal gas equation:
1. Write down all the information that you know about the gas.
2. If necessary, convert the known values to SI units.
3. Choose a relevant gas law equation that will allow you to calculate the unknown variable.
4. Substitute the known values into the equation. Calculate the unknown variable.
Remember that the general gas equation only applies if the molar quantity of the gas is fixed. For example, if a gas is mixed with
another gas, you may have to apply the equation separately for individual gases.
Example
Let’s imagine that at the beginning of a journey a truck tire has a volume of 30,000 cm3 and an internal pressure of 170 kPa. The
temperature of the tire is 16∘C. By the end of the trip, the volume of the tire has increased to 32,000 cm3 and the temperature of the
air inside the tire is 40∘C. What is the tire pressure at the end of the journey?
Tire Pressure: Tire pressure may change significantly during the operation of the vehicle. This is mostly due to the temperature
change of the air in tires.
Solution:
Step 1. Write down all the information that you know about the gas: P1= 170 kPa and P2 is unknown. V1 = 30,000 cm3 and V2 =
32,000 cm3. T1 = 16∘C and T2 = 40∘C.
Step 2. Convert the known values to SI units if necessary: Here, temperature must be converted into Kelvin. Therefore, T1 = 16 +
273 = 289 K, T2 = 40 + 273 = 313 K
Step 3. Choose a relevant gas law equation that will allow you to calculate the unknown variable: We can use the general gas
P1 V1 P2 V2
equation to solve this problem: =
T1
. T2
P1 ×V1 ×T2
Therefore, P 2 =
T1 ×V2
.
Step 4. Substitute the known values into the equation. Calculate the unknown variable:
12.4.5 https://phys.libretexts.org/@go/page/14511
170 × 30, 000 × 313
P2 = = 173 kPa (12.4.16)
289 × 32, 000
The pressure of the tire at the end of the journey is 173 kPa.
Note that in Step 2 we did not bother to convert the volume values to m3. In Step 4, pressure appears both in the numerator and
denominator. In this case the conversion was not necessary.
Avogador’s Number
The number of molecules in a mole is called Avogadro’s number (NA)—defined as 6.02x 1023 mol-1.
learning objectives
Explain relationship between Avogadro’s number and mole.
When measuring the amount of substance, it is sometimes easier to work with a unit other than the number of molecules. A mole
(abbreviated mol) is a base unit in the International System of Units (SI). It is defined as any substance containing as many atoms
or molecules as there are in exactly 12 grams (0.012 kg) of carbon-12. The actual number of atoms or molecules in one mole is
called Avogadro’s constant (NA), in recognition of Italian scientist Amedeo Avogadro.
Amadeo Avogadro: Amedeo Avogadro (1776–1856). He established the relationship between the masses of the same volume of
different gases (at the same temperature and pressure) corresponds to the relationship between their respective molecular weights.
Avogadro’s number (N) refers to the number of molecules in one gram-molecule of oxygen. This indicates an amount of substance
as opposed to an independent dimension of measurement. In 1811 Amedeo Avogadro first proposed that the volume of a gas (at a
given pressure and temperature) is proportional to the number of atoms or molecules, regardless of the nature of the gas (i.e., this
number is universal and independent of the type of gas). In 1926, Jean Perrin won the Nobel Prize in Physics, largely for his work
in determining the Avogadro constant (by several different methods). The value of Avogadro’s constant, NA , has been found to
equal 6.02×1023mol−1.
Role in Science
Avogadro’s constant is a scaling factor between macroscopic and microscopic (atomic scale) observations of nature. As such, it
provides the relation between other physical constants and properties. For example, it establishes a relationship between the gas
constant R and the Boltzmann constant k,
12.4.6 https://phys.libretexts.org/@go/page/14511
−1 −1
R = kNA = 8.314472(15)J mo l K ; (12.4.17)
Measuring NA
The determination of NA is crucial to the calculation of an atom’s mass, since the latter is obtained by dividing the mass of a mole
of the gas by Avogadro’s constant. In his study on Brownian motion in 1905, Albert Einstein proposed that this constant could be
determined based on the quantities observable in Brownian motion. Subsequently, Einstein’s idea was verified, leading to the first
determination of NA in 1908 through the experimental work of Jean Baptiste Perrin.
Absolute Temperature
Absolute temperature is the most commoly used thermodyanmic temperature unit and is the standard unit of temperature.
learning objectives
Describe relationship between absolute temperature and kinetic energy.
Thermodynamic temperature is the absolute measure of temperature. It is one of the principal parameters of thermodynamics and
kinetic theory of gases. Thermodynamic temperature is an “absolute” scale because it is the measure of the fundamental property
underlying temperature: its null or zero point (“absolute zero”) is the temperature at which the particle constituents of matter have
minimal motion and cannot become any colder. That is, they have minimal motion, retaining only quantum mechanical motion, as
diagramed in.
Graph of Pressure Versus Temperature: Graph of pressure versus temperature for various gases kept at a constant volume. Note
that all of the graphs extrapolate to zero pressure at the same temperature
At its simplest, “temperature” arises from the kinetic energy of the random motions of matter’s particle constituents such as
molecules or atoms, as seen in. Therefore, it is reasonable to choose absolute zero, where all classical motion ceases, as the
reference point (T=0) of our temperature system. By using the absolute temperature scale (Kelvin system), which is the most
commonly used thermodynamic temperature, we have shown that the average translational kinetic energy (KE) of a particle in a
gas has a simple relationship to the temperature:
12.4.7 https://phys.libretexts.org/@go/page/14511
Translational Motion of Helium: Real gases do not always behave according to the ideal model under certain conditions, such as
high pressure. Here, the size of helium atoms relative to their spacing is shown to scale under 1950 atmospheres of pressure.
3
¯
KE = kT. (12.4.19)
2
Note that this equation would not look this elegant if the Fahrenheit scale were used instead.
Key Points
Ideal gas law was derived empirically by combining Boyle’s law and Charles’ law.
Although the empirical derivation of the equation does not consider microscopic details, the ideal gas law can be derived from
first principles in the classical thermodynamics.
Pressure and volume of a gas can be related to the average velocity of molecues: PV = Nmv¯ .1
3
2
Isothermal processes typically occur when a system is in contact with an outside thermal reservoir ( heat bath), and the change
occurs slowly enough to allow the system to adjust continually to the temperature of the reservoir through heat exchange.
For an ideal gas, from the ideal gas law PV = NkT, PV remains constant through an isothermal process. A curve in a P-V
diagram generated by the equation PV = const is called an isotherm.
For an isothermal, reversible process, the work done by the gas is equal to the area under the relevant pressure -volume
VB
isotherm. It is given as W → B = NkT ln
A
VA
.
Gases can expand or contract under a certain constraint. Depending on the constraint, the final state of the gas may change.
The heat transferred to the system does work but also changes the internal energy of the system. In an isobaric process for a
monatomic gas, heat and the temperature change satisfy the following equation: Q = NkΔT .
5
Write down all the information that you know about the gas and convert the known values to SI units if necessary.
Choose a relevant gas law equation that will allow you to calculate the unknown variable, and substitute the known values into
the equation. Then calculate the unknown variable.
The general gas equation only applies if the molar quantity of the gas is fixed.
Avogadro hypothesized that equal volumes of gas, at the same pressure and temperature, contain equal numbers of molecules,
regardless of the type of gas.
Avogadro’s constant is a scaling factor between macroscopic and microscopic (atomic scale) observations of nature. It provides
the relation between other physical constants and properties.
Albert Einstein proposed that Avogadro’s number could be determined based on the quantities observable in Brownian motion.
NA was measured for the first time by Jean Baptiste Perrin in 1908.
Temperature arises from the kinetic energy of the random motions of matter ‘s particle constituents such as molecules or atoms.
Therefore, it is reasonable to choose absolute zero, where all classical motion ceases, as the reference point.
12.4.8 https://phys.libretexts.org/@go/page/14511
By international agreement, the unit kelvin and its scale are defined by two points: absolute zero and the triple point of the
standardized water.
At absolute zero, the particle constituents of matter have minimal motion and cannot become any colder. They retain minimal,
quantum mechanical motion.
Key Terms
mole: In the International System of Units, the base unit of amount of substance; the amount of substance of a system which
contains as many elementary entities as there are atoms in 12 g of carbon-12. Symbol: mol.
ideal gas: A hypothetical gas whose molecules exhibit no interaction and undergo elastic collision with each other and with the
walls of the container.
Avogadro’s number: the number of constituent particles (usually atoms or molecules) in one mole of a given substance. It has
dimensions of reciprocal mol and its value is equal to 6.02214129·1023 mol-1
adiabatic: Occurring without gain or loss of heat.
internal energy: The sum of all energy present in the system, including kinetic and potential energy; equivalently, the energy
needed to create a system, excluding the energy necessary to displace its surroundings.
the first law of thermodynamics: A version of the law of energy conservation: the change in the internal energy of a closed
system is equal to the amount of heat supplied to the system, minus the amount of work done by the system on its surroundings.
specific heat: The ratio of the amount of heat needed to raise the temperature of a unit mass of substance by a unit degree to the
amount of heat needed to raise that of the same mass of water by the same amount.
SI units: International System of Units (abbreviated SI from French: Le Système international d’unités). It is the modern form
of the metric system.
gas constant: A universal constant, R, that appears in the ideal gas law, (PV = nRT), derived from two fundamental constants,
the Boltzman constant and Avogadro’s number, (R = NAk).
Faraday constant: The magnitude of electric charge per mole of electrons.
Brownian motion: Random motion of particles suspended in a fluid, arising from those particles being struck by individual
molecules of the fluid.
absolute zero: The coldest possible temperature: zero on the Kelvin scale and approximately -273.15°C and -459.67°F. The
total absence of heat; the temperature at which motion of all molecules would cease.
International System of Units: (SI): The standard set of basic units of measurement used in scientific literature worldwide.
Vienna Standard Mean Ocean Water: A standard defining a standardized isotopic composition of water.
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12.4.11 https://phys.libretexts.org/@go/page/14511
12.5: Kinetic Theory
learning objectives
Express the relationship between the pressure and the average kinetic energy of gas molecules in the form of equation
In Newtonian mechanics, if pressure is the force divided by the area on which the force is exerted, then what is the origin of
pressure in a gas? What forces create the pressure? We can gain a better understanding of pressure (and temperature as well) from
the kinetic theory of gases, which assumes that atoms and molecules are in continuous random motion.
Translational Motion of Helium: Real gases do not always behave according to the ideal model under certain conditions, such as
high pressure. Here, the size of helium atoms relative to their spacing is shown to scale under 1950 atmospheres of pressure.
Δp = pi,x − pf ,x = pi,x − (−pi,x ) (12.5.1)
vx
,
(where L is the distance between opposite walls). The force due to this particle is:
2
Δp mvx
F = = . (12.5.3)
Δt L
where the bar denotes an average over the N particles. Since the assumption is that the particles move in random directions, if we
divide the velocity vectors of all particles in three mutually perpendicular directions, the average value of the squared velocity
along each direction must be same. (This does not mean that each particle always travel in 45 degrees to the coordinate axes. )
12.5.1 https://phys.libretexts.org/@go/page/14512
Pressure: Pressure arises from the force exerted by molecules or atoms impacting on the walls of a container.
¯2 ¯2
This gives v¯
2
x =
v
3
. We can rewrite the force as F = Nmv
3L
.
This force is exerted on an area L2. Therefore the pressure of the gas is:
¯2 ¯2
F Nmv nmv
P = = = , (12.5.5)
2
L 3V 3
where V=L3 is the volume of the box. The fraction n=N/V is the number density of the gas. This is a first non-trivial result of the
kinetic theory because it relates pressure (a macroscopic property) to the average (translational) kinetic energy per molecule which
is a microscopic property.
learning objectives
Describe the shape and temperature dependence of the Maxwell-Boltzmann distribution curve
The motion of molecules in a gas is random in magnitude and direction for individual molecules, but a gas of many molecules has a
predictable distribution of molecular speeds, known as the Maxwell-Boltzmann distribution (illustrated in ). The distribution has a
long tail because some molecules may go several times the rms speed. The most probable speed vp (at the peak of the curve) is less
than the rms speed vrms. As shown in, the curve is shifted to higher speeds at higher temperatures, with a broader range of speeds.
Maxwell-Boltzmann Distribution at Higher Temperatures: The Maxwell-Boltzmann distribution is shifted to higher speeds and
is broadened at higher temperatures.
12.5.2 https://phys.libretexts.org/@go/page/14512
Maxwell-Boltzmann Distribution: The Maxwell-Boltzmann distribution of molecular speeds in an ideal gas. The most likely
speed v_p is less than the rms speed v_rms. Although very high speeds are possible, only a tiny fraction of the molecules have
speeds that are an order of magnitude greater than v_rms.
Maxwell-Boltzmann Distribution
Maxwell-Boltzmann distribution is a probability distribution. It applies to ideal gases close to thermodynamic equilibrium, and is
given as the following equation:
2 2 2
m m(vx + vy + vz )
3/2
fv (vx , vy , vz ) = ( ) exp[− ], (12.5.6)
2πkT 2kT
where fv is the velocity probability density function. (Derivation of the formula goes beyond the scope of introductory physics. )
The formula calculates the probability of finding a particle with velocity in the infinitesimal element [dvx, dvy, dvz] about velocity v
= [vx, vy, vz] is:
fv (vx , vy , vz )dvx dvy dvz . (12.5.7)
It can also be shown that the Maxwell–Boltzmann velocity distribution for the vector velocity [vx, vy, vz] is the product of the
distributions for each of the three directions:
This makes sense because particles are moving randomly, meaning that each component of the velocity should be independent.
where θ and ϕ are the “course” (azimuth of the velocity vector) and “path angle” (elevation angle of the velocity vector).
Integration of the normal probability density function of the velocity, above, over the course (from 0 to 2π) and path angle (from 0
to π), with substitution of the speed for the sum of the squares of the vector components, yields the following probability density
function (known simply as the Maxwell distribution):
−−−−−− 2
f (v) = √(
m
2πkT
3 2
) 4π v exp(
−mv
2kT
) for speed v.
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Temperature
Temperature is directly proportional to the average translational kinetic energy of molecules in an ideal gas.
learning objectives
Describe relationship between temperature and energy of molecules in an ideal gas
Intuitively, hotter air suggests faster movement of air molecules. In this atom, we will derive an equation relating the temperature
of a gas (a macroscopic quantity) to the average kinetic energy of individual molecules (a microscopic quantity). This is a basic and
extremely important relationship in the kinetic theory of gases.
Microscopic View
We assume that a molecule is small compared with the separation of molecules in the gas (confined in a three dimensional
container), and that its interaction with other molecules can be ignored. Also, we assume elastic collisions when molecules hit the
wall of the container, as illustrated in.
Elastic Collisions When Molecules Hit the Wall of the Container: Gas in a box exerts an outward pressure on its walls. A
molecule colliding with a rigid wall has the direction of its velocity and momentum in the x-direction reversed. This direction is
perpendicular to the wall. The components of its velocity momentum in the y- and z-directions are not changed, which means there
is no force parallel to the wall.
We have seen in the Atom on “Origin of Pressure ” that, for an ideal gas under our assumptions:
¯2
Nmv
P = , (12.5.12)
3V
where P is the pressure, N is the number of molecules, m is the mass of the molecule, v is the speed of molecules, and V is the
volume of the gas. From the equation, we get:
1 ¯2
PV = Nmv (Eq. 1). (12.5.13)
3
What can we learn from this atomic and molecular version of the ideal gas law ? We can derive a relationship between temperature
and the average translational kinetic energy of molecules in a gas. Recall the macroscopic expression of the ideal gas law:
where N is the number of molecules, T is the temperature of the gas, and k is the Boltzmann constant.
Equating the right hand sides of the macroscopic and microscopic versions of the ideal gas law (Eq. 1 & 2) gives:
1 ¯2
mv = kT. (12.5.15)
3
12.5.4 https://phys.libretexts.org/@go/page/14512
Thermal Energy
Note that the average kinetic energy (KE) of a molecule in the gas is:
1 ¯2
mv . (12.5.16)
2
RMS Speed
Eq. 3 is a molecular interpretation of temperature. It has been found to be valid for gases and reasonably accurate in liquids and
solids. It is another definition of temperature based on an expression of the molecular energy. It is sometimes useful to know the
average speed of molecules in a gas in terms of temperature:
−−−−
¯2 3kT
v = vrms = √ , (12.5.18)
m
learning objectives
Determine the number of degrees of freedom and calculate the internal energy for an ideal gas molecule
In thermodynamics, internal energy is the total energy contained by a thermodynamic system. Internal energy has two major
components: kinetic energy and potential energy. The kinetic energy is due to the motion of the system’s particles (e.g.,
translations, rotations, vibrations). In ideal gases, there is no inter-particle interaction. Therefore, we will disregard potential energy
and only focus on the kinetic energy contribution to the internal energy.
Monatomic Gases
A monatomic gas is one in which atoms are not bound to each other. Noble gases (He, Ne, etc.) are typical examples. A helium
balloon is shown in the following figure. In this case, the kinetic energy consists only of the translational energy of the individual
atoms. Monoatomic particles do not vibrate, and their rotational energy can be neglected because atomic moment of inertia is so
small. Also, they are not electronically excited to higher energies except at very high temperatures. Therefore, practical internal
energy changes in an ideal gas may be described solely by changes in its translational kinetic energy.
Helium Blimp: Helium, like other noble gases, is a monatomic gas, which often can be described by the ideal gas law. It is the gas
of choice to fill airships such as the Goodyear blimp.
The average kinetic energy (KE) of a particle in an ideal gas is given as:
1 ¯2 3
¯
KE = mv = kT, (12.5.19)
2 2
12.5.5 https://phys.libretexts.org/@go/page/14512
where k is the Boltzmann’s constant. (See the Atom on “Temperature” in kinetic theory. ) With N atoms in the gas, its total internal
energy U is given as:
3
U = NkT (12.5.20)
2
where N is the number of atoms in the gas. Note that there are three degrees of freedom in monatomic gases: translation in x, y and
z directions.
Since atomic motion is random (and therefore isotropic), each degrees of freedom contribute 1
2
kT per atom to the internal energy.
Diatomic gases
A diatomic molecule (H2, O2, N2, etc.) has 5 degrees of freedom (3 for translation in x, y and z directions, and 2 for rotation).
Therefore, the internal energy for diatomic gases is U = NkT. 5
Key Points
We can gain a better understanding of pressure (and temperature as well) from the kinetic theory of gases, which assumes that
atoms and molecules are in continuous random motion.
Pressure, a macroscopic property, can be related to the average (translational) kinetic energy per molecule which is a
¯
microscopic property by P = nmv2
3
.
Since the assumption is that the particles move in random directions, the average value of velocity squared along each direction
must be same. This gives: v¯ = v¯ = v¯ = v¯ /3 .
2
x
2
y
2
z
2
The Maxwell-Boltzmann distribution has a long tail, and the most probable speed vp is less than the rms speed vrms. The
distribution curve is shifted to higher speeds at higher temperatures, with a broader range of speeds.
2 2 2
m( vx +vy +vz )
Maxwell-Boltzmann distribution is given as follows: f v (vx , vy , vz ) = (
m
2πkT
)
3/2
exp[−
2kT
] . It is a product of three
independent 1D Maxwell-Boltzmann distributions.
−−−−−− 2
2πkT
3
) 4π v
2
exp(
−mv
2kT
) . This is simply called Maxwell distribution.
The average translational kinetic energy of a molecule is equivalent to kT and is called thermal energy.
3
In kinematic theory of gases, macroscopic quantities (such as press and temperature) are explained by considering microscopic
(random) motion of molecules.
−−−
The rms speed of molecules in a gas is given as √ 3kT
m
.
In ideal gases, there is no inter-particle interaction. Therefore, only the kinetic energy contribute to the internal energy.
Each degrees of freedom contribute kT per atom to the internal energy.
1
For monatomic ideal gases with N atoms, its total internal energy U is given as U = NkT . For diatomic gases, U =
3
2
5
2
NkT .
Key Terms
kinetic theory of gases: The kinetic theory of gases describes a gas as a large number of small particles (atoms or molecules),
all of which are in constant, random motion.
Newtonian mechanics: Early classical mechanics as propounded by Isaac Newton, especially that based on his laws of motion
and theory of gravity.
rms: Root mean square: a statistical measure of the magnitude of a varying quantity.
ideal gas: A hypothetical gas whose molecules exhibit no interaction and undergo elastic collision with each other and with the
walls of the container.
Boltzmann’s constant: The physical constant relating energy at the particle level with temperature observed at the bulk level. It
is the gas constant R divided by Avogadro’s number, NA.
moment of inertia: A measure of a body’s resistance to a change in its angular rotation velocity
noble gas: Any of the elements of group 18 of the periodic table, being monatomic and (with very limited exceptions) inert.
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12.6: Phase Changes
learning objectives
Describe behavior of the medium during a phase transition
A phase of a thermodynamic system and the states of matter have uniform physical properties. During a phase transition of a given
medium certain properties of the medium change, often discontinuously, as a result of some external condition, such as temperature
or pressure. For example, a liquid may become gas upon heating to the boiling point, resulting in an abrupt change in volume. The
measurement of the external conditions at which the transformation occurs is termed the phase transition. The term is most
commonly used to describe transitions between solid, liquid and gaseous states of matter and, in rare cases, plasma.
As an example, if you boil water, it never goes above 100 degrees Celsius. Only after it has completely evaporated will it get any
hotter. This is because once water reaches the boiling point, extra energy is used to change the state of matter and increase the
potential energy instead of the kinetic energy. The opposite happens when water freezes. To boil or melt one mole of a substance, a
certain amount of energy is required. These amounts of energy are the molar heat of vaporization and molar heat of fusion. If that
amount of energy is added to a mole of that substance at boiling or freezing point, all of it will melt or boil, but the temperature
won’t change.
Temperature increases linearly with heat, until the melting point. But the heat added does not change the temperature; that heat
energy is instead used to break intermolecular bonds and convert ice into water. At this point, there is a mixture of both ice and
water. Once all ice has been melted, the temperature again rises linearly with heat added. At the boiling point, temperature no
longer rises with heat added because the energy is once again being used to break intermolecular bonds. Once all water has been
boiled to steam, the temperature will continue to rise linearly as heat is added.
Temperature vs. Heat: This graph shows the temperature of ice as heat is added.
The plots of pressure versus temperatures provide considerable insight into thermal properties of substances. There are well-
defined regions on these graphs that correspond to various phases of matter, so PT graphs are called phase diagrams. Using the
graph, if you know the pressure and temperature you can determine the phase of water. The solid lines—boundaries between
phases—indicate temperatures and pressures at which the phases coexist (that is, they exist together in ratios, depending on
pressure and temperature). For example, the boiling point of water is 100º C at 1.00 atm. As the pressure increases, the boiling
temperature rises steadily to 374º C at a pressure of 218 atm. A pressure cooker (or even a covered pot) will cook food faster
because the water can exist as a liquid at temperatures greater than 100º C without all boiling away. The curve ends at a point called
the critical point, because at higher temperatures the liquid phase does not exist at any pressure. The critical temperature for oxygen
is -118ºC, so oxygen cannot be liquefied above this temperature.
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Phase Diagram of Water: In this typical phase diagram of water, the green lines mark the freezing point, and the blue line marks
the boiling point, showing how they vary with pressure. The dotted line illustrates the anomalous behavior of water. Note that water
changes states based on the pressure and temperature.
learning objectives
Explain why water boils at 100 °C
Overview
The term relative humidity refers to how much water vapor is in the air compared with the maximum possible. At its maximum,
denoted as saturation, the relative humidity is 100%, and evaporation is inhibited. The amount of water vapor the air can hold
depends on its temperature. For example, relative humidity rises in the evening, as air temperature declines, sometimes reaching the
dew point. At the dew point temperature, relative humidity is 100%, and fog may result from the condensation of water droplets if
they are small enough to stay in suspension. Conversely, if one wished to dry something, it is more effective to blow hot air over it
rather than cold air, because, among other things, hot air can hold more water vapor.
Evaporation
The capacity of air to hold water vapor is based on vapor pressure of water. The liquid and solid phases are continuously giving off
vapor because some of the molecules have high enough speeds to enter the gas phase, a process called evaporation; see (a). For the
molecules to evaporate, they must be located near the surface, be moving in the proper direction, and have sufficient kinetic energy
to overcome liquid-phase intermolecular forces. When only a small proportion of the molecules meet these criteria, the rate of
evaporation is low. Since the kinetic energy of a molecule is proportional to its temperature, evaporation proceeds more quickly at
higher temperatures.
If a lid is placed over the container, as in (b), evaporation continues, increasing the pressure, until sufficient vapor has built up for
condensation to balance evaporation. Then equilibrium has been achieved, and the vapor pressure is equal to the partial pressure of
water in the container. Vapor pressure increases with temperature because molecular speeds are higher as temperature increases.
As the faster-moving molecules escape, the remaining molecules have lower average kinetic energy, and the temperature of the
liquid decreases. This phenomenon is also called evaporative cooling. This is why evaporating sweat cools the human body.
Evaporation also tends to proceed more quickly with higher flow rates between the gaseous and liquid phase and in liquids with
higher vapor pressure. For example, laundry on a clothes line will dry (by evaporation) more rapidly on a windy day than on a still
day.
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starts out at the bottom of the container at 20ºC, it contains water vapor (about 2.30%). The pressure inside the bubble is fixed at
1.00 atm (we ignore the slight pressure exerted by the water around it). As the temperature rises, the amount of air in the bubble
stays the same, but the water vapor increases; the bubble expands to keep the pressure at 1.00 atm. At 100ºC, water vapor enters the
bubble continuously since the partial pressure of water is equal to 1.00 atm in equilibrium. It cannot reach this pressure, however,
since the bubble also contains air and total pressure is 1.00 atm. The bubble grows in size and thereby increases the buoyant force.
The bubble breaks away and rises rapidly to the surface, resulting in boiling. (See. )
Close-up of the Boiling Process: (a) An air bubble in water starts out saturated with water vapor at 20ºC. (b) As the temperature
rises, water vapor enters the bubble because its vapor pressure increases. The bubble expands to keep its pressure at 1.00 atm. (c)
At 100ºC, water vapor enters the bubble continuously because water’s vapor pressure exceeds its partial pressure in the bubble,
which must be less than 1.00 atm. The bubble grows and rises to the surface.
Key Points
The term is most commonly used to describe transitions between solid, liquid and gaseous states of matter and, in rare cases,
plasma.
Once water reaches the boiling point, extra energy is used to change the state of matter and increase the potential energy instead
of the kinetic energy.
Plots of pressure versus temperatures, an example of a phase diagram, provide considerable insight into thermal properties of
substances.
Relative humidity is the fraction of water vapor in a gas compared to the saturation value.
Since the kinetic energy of a molecule is proportional to its temperature, evaporation proceeds more quickly at higher
temperatures.
Vapor pressure increases with temperature because molecular speeds are higher as temperature increases.
Water boils at 100 °C because the vapor pressure exceeds atmospheric pressure at this temperature.
Key Terms
intermolecular: from one molecule to another; between molecules
plasma: a state of matter consisting of partially ionized gas
thermodynamic: Relating to the conversion of heat into other forms of energy.
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
vapor pressure: The pressure that a vapor exerts, or the partial pressure if it is mixed with other gases.
humidity: The amount of water vapor in the air.
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12.7: The Zeroth Law of Thermodynamics
learning objectives
Identify major implications of the Zeroth Law of Thermodynamics
Zeroth Law of Thermodynamics: The double arrow represents thermal equilibrium between systems. If systems A and C are in
equilibrium, and systems A and B are in equilibrium, then systems B and C are in equilibrium. The systems A, B, and C are at the
same temperature.
Key Points
Assuming A, B, and C are three systems, if A and C are in thermal equilibrium, and A and B are in thermal equilibrium, then B
and C are in thermal equilibrium.
Two systems are in thermal equilibrium if they could transfer heat between each other, but do not.
The Zeroth Law of Thermodynamics implies that temperature is a quantity worth measuring.
Key Terms
thermal equilibrium: Two systems are in thermal equilibrium if they could transfer heat between each other, but don’t.
zeroth law of thermodynamics: Let A, B and C be three systems. If A and C are in thermal equilibrium, and A and B are in
thermal equilibrium, then B and C are in thermal equilibrium.
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12.8: Thermal Stresses
learning objectives
Formulate relationship between thermal stress and thermal expansion
Thermal Expansion
Thermal expansion is the change in size or volume of a given mass with temperature. The expansion of alcohol in a thermometer is
one of many commonly encountered examples of this. Hot air rises because its volume increases, which causes the hot air’s density
to be smaller than the density of surrounding air, causing a buoyant (upward) force on the hot air. The same happens in all liquids
and gases, driving natural heat transfer upward in homes, oceans, and weather systems. Solids also undergo thermal expansion.
Railroad tracks and bridges, for example, have expansion joints to allow them to freely expand and contract with temperature
changes.
Thermal Expansion Joints: Thermal expansion joints like these in the Auckland Harbour Bridge in New Zealand allow bridges to
change length without buckling. (credit: Ingolfson, Wikimedia Commons)
What are the basic properties of thermal expansion? First, thermal expansion is clearly related to temperature change. The greater
the temperature change, the more a bimetallic strip will bend. Second, it depends on the material. In a thermometer, for example,
the expansion of alcohol is much greater than the expansion of the glass containing it.
What is the underlying cause of thermal expansion? An increase in temperature implies an increase in the kinetic energy of the
individual atoms. In a solid, unlike in a gas, the atoms or molecules are closely packed together, but their kinetic energy (in the
form of small, rapid vibrations) pushes neighboring atoms or molecules apart from each other. This neighbor-to-neighbor pushing
results in a slightly greater distance, on average, between neighbors, and adds up to a larger size for the whole body. For most
substances under ordinary conditions, there is no preferred direction, and an increase in temperature will increase the solid’s size by
a certain fraction in each dimension.
To be more quantitative, the change in length ΔL is proportional to length L. The dependence of thermal expansion on temperature,
substance, and length is summarized in the equation
ΔL = αLΔT (12.8.1)
where ΔL is the change in length L, ΔT is the change in temperature, and α is the coefficient of linear expansion, which varies
slightly with temperature.
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Thermal Stress
Thermal stress is created by thermal expansion or contraction. Thermal stress can be destructive, such as when expanding gasoline
ruptures a tank. It can also be useful, for example, when two parts are joined together by heating one in manufacturing, then
slipping it over the other and allowing the combination to cool. Thermal stress can explain many phenomena, such as the
weathering of rocks and pavement by the expansion of ice when it freezes.
Forces and pressures created by thermal stress can be quite large. Railroad tracks and roadways can buckle on hot days if they lack
sufficient expansion joints. Power lines sag more in the summer than in the winter, and will snap in cold weather if there is
insufficient slack. Cracks open and close in plaster walls as a house warms and cools. Glass cooking pans will crack if cooled
rapidly or unevenly, because of differential contraction and the stresses it creates. (Pyrex® is less susceptible because of its small
coefficient of thermal expansion. ) Nuclear reactor pressure vessels are threatened by overly rapid cooling, and although none have
failed, several have been cooled faster than considered desirable. Biological cells are ruptured when foods are frozen, detracting
from their taste. Repeated thawing and freezing accentuates the damage. Even the oceans can be affected. A significant portion of
the rise in sea level that is resulting from global warming is due to the thermal expansion of sea water.
Metal is regularly used in the human body for hip and knee implants. Most implants need to be replaced over time because, among
other things, metal does not bond with bone. Researchers are trying to find better metal coatings that would allow metal-to-bone
bonding. One challenge is to find a coating that has an expansion coefficient similar to that of metal. If the expansion coefficients
are too different, the thermal stresses during the manufacturing process lead to cracks at the coating-metal interface.
Another example of thermal stress is found in the mouth. Dental fillings can expand differently from tooth enamel. It can give pain
when eating ice cream or having a hot drink. Cracks might occur in the filling. Metal fillings (gold, silver, etc.) are being replaced
by composite fillings (porcelain), which have smaller coefficients of expansion, and are closer to those of teeth.
Key Points
Thermal expansion is the change in size or volume of a given mass with changing temperature.
An increase in temperature implies an increase in the kinetic energy of the individual atoms, which will increase a solid’s size
by a certain fraction in each dimension.
Thermal stress is created when thermal expansion is constrained.
Key Terms
stress: The internal distribution of force per unit area (pressure) within a body reacting to applied forces which causes strain or
deformation and is typically symbolized by σ.
differential: A qualitative or quantitative difference between similar or comparable things.
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12.9: Diffusion
learning objectives
Discuss the process and results of diffusion, identifying factors that affect its rate
Diffusion is the movement of particles move from an area of high concentration to an area of low concentration until equilibrium is
reached. A distinguishing feature of diffusion is that it results in mixing or mass transport without requiring bulk motion. Thus,
diffusion should not be confused with convection or advection, which are other transport mechanisms that use bulk motion to move
particles from one place to another.
Diffusion: Particles moving from areas of high concentration to areas of low concentration.
Molecular diffusion, often called simply diffusion, is the thermal motion of all (liquid or gas) particles at temperatures above
absolute zero. The rate of this movement is a function of temperature, viscosity of the fluid and the size (mass) of the particles.
Diffusion explains the net flux of molecules from a region of higher concentration to one of lower concentration. However,
diffusion can still occur in the absence of a concentration gradient.
The result of diffusion is a gradual mixing of material. In a phase with uniform temperature, absent external net forces acting on the
particles, the diffusion process will eventually result in complete mixing.
Key Points
Molecular diffusion, often called simply diffusion, is the thermal motion of all (liquid or gas) particles at temperatures above
absolute zero.
The result of diffusion is a gradual mixing of material. In a phase with uniform temperature, absent external net forces acting on
the particles, the diffusion process will eventually result in complete mixing.
Diffusion can also occur in the absence of a concentration gradient — equilibrium particles are still moving around their
container.
Key Terms
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
diffusion: Diffusion is the movement of particles from regions of high concentration toward regions of lower concentration.
concentration: The proportion of a substance in a mixture.
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CHAPTER OVERVIEW
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1
13.1: Introduction
learning objectives
Distinguish three modes of heat transfer
Examples of Heat Transfer: (a) The chilling effect of a clear breezy night is produced by the wind and by radiative heat transfer to
cold outer space. (b) There was once great controversy about the Earth’s age, but it is now generally accepted to be about 4.5
billion years old. Much of the debate is centered on the Earth’s molten interior. According to our understanding of heat transfer, if
the Earth is really that old, its center should have cooled off long ago. The discovery of radioactivity in rocks revealed the source of
energy that keeps the Earth’s interior molten, despite heat transfer to the surface, and from there to cold outer space.
Definitions
Scottish physicist James Clerk Maxwell, in his 1871 classic Theory of Heat, was one of many who began to build on the already
established idea that heat was something to do with matter in motion. This was the same idea put forwards by Sir Benjamin
Thompson in 1798, who said he was only following on from the work of many others. One of Maxwell’s recommended books was
Heat as a Mode of Motion by John Tyndall. Maxwell outlined four stipulations for the definition of heat:
It is something which may be transferred from one body to another.
It is a measurable quantity, and thus treated mathematically.
It cannot be treated as a substance, because it may be transformed into something that is not a substance, such as mechanical
work.
Heat is one of the forms of energy.
In the following sections, we will define heat more rigorously, paying particular attention to how it can be measured and quantified.
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Heat Transfer Methods
After defining and quantifying heat transfer and its effects on physical systems, we will discuss the methods by which heat is
transferred. So many processes involve heat transfer, so that it is hard to imagine a situation where no heat transfer occurs. Yet
every process involving heat transfer takes place by only three methods:
1. Conduction is heat transfer through stationary matter by physical contact. ( The matter is stationary on a macroscopic scale—we
know there is thermal motion of the atoms and molecules at any temperature above absolute zero. ) Heat transferred between
the electric burner of a stove and the bottom of a pan is transferred by conduction.
2. Convection is the heat transfer by the macroscopic movement of a fluid. This type of transfer takes place in a forced-air furnace
and in weather systems, for example.
3. Heat transfer by radiation occurs when microwaves, infrared radiation, visible light, or another form of electromagnetic
radiation is emitted or absorbed. An obvious example is the warming of the Earth by the Sun. A less obvious example is
thermal radiation from the human body.
learning objectives
Identify SI and common units of heat
Heat Transfer and Equilibrium: (a) The soft drink and the ice have different temperatures, T1 and T2, and are not in thermal
equilibrium. (b) When the soft drink and ice are allowed to interact, energy is transferred until they reach the same temperature T,
achieving equilibrium. Heat transfer occurs due to the difference in temperatures. In fact, since the soft drink and ice are both in
contact with the surrounding air and bench, the equilibrium temperature will be the same for both.
Where Is the Most Heat Lost?: Use movable thermometers to discover where a house has poor insulation.
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AP Physics 2 - Heat Transfer
Units
Owing to the fact that heat is a form of energy, it has the SI unit of joule (J). The calorie (cal) is a common unit of energy, defined
as the energy needed to change the temperature of 1.00 g of water by 1.00ºC —specifically, between 14.5ºC and 15.5ºC, since there
is a slight temperature dependence. Another common unit of heat is the kilocalorie (kcal), which is the energy needed to change the
temperature of 1.00 kg of water by 1.00ºC. Since mass is often specified in kilograms, kilocalorie is commonly used. Food calories
(given the notation Cal, and sometimes called “big calorie”) are actually kilocalories (1kilocalorie=1000 calories), a fact not easily
determined from package labeling in the United States, but more common in Europe and elsewhere. In some engineering fields, the
British Thermal Unit (BTU), equal to about 1.055 kilo-joules, is widely used.
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Figure 1 Equivalence of Heat and Work: Schematic depiction of Joule’s experiment that established the equivalence of heat and
work
The total amount of energy transferred as heat is conventionally written as Q for algebraic purposes. Heat released by a system into
its surroundings is by convention a negative quantity (Q < 0); when a system absorbs heat from its surroundings, it is positive (Q >
0).
Internal Energy
The internal energy of a system is the sum of all kinetic and potential energy in a system.
learning objectives
Express the internal energy in terms of kinetic and potential energy
Internal Energy
James Joule showed that both heat and work can produce the same change in the internal energy of a substance, establishing the
principle of the mechanical equivalence of heat. Heat is emphatically a quantity that solely describes energy being transferred. It
makes no sense to speak of the total ‘heat’ an object or system contains. However, a system does contain a quantifiable amount of
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energy called the internal energy of a system. The internal energy of a system is the quantity that changes with the addition or
subtraction of work or heat. It is closely related to temperature.
Definition
The internal energy is the energy required to create a system, excluding the energy necessary to displace its surroundings. Internal
energy has two components: kinetic energy and potential energy. The kinetic energy consists of all the energy involving the
motions of the particles constituting the system, including translation, vibration, and rotation. The potential energy is associated
with the static constituents of matter, static electric energy of atoms within molecules or crystals, and the energy from chemical
bonds. The equation describing the total internal energy of a system is then:
U = Ukinetic + Upotential. (13.1.1)
We can also think of the internal energy as the sum of all the energy states of each component in the system:
U = ∑ Ei . (13.1.2)
i
At any finite temperature, kinetic and potential energies are constantly converted into each other, but the total energy remains
constant in an isolated system. The kinetic energy portion of internal energy gives rise to the temperature of the system. We can use
statistical mechanics to relate the (somewhat) random motions of particles in a system to the mean kinetic energy of the ensemble
of particles, and thus the empirically measurable quantity expressed as temperature.
We can see that internal energy is an extensive property: it depends on the size of the system or on the amount of substance it
contains.
In most cases, we are not concerned with the total amount of internal energy in the system, as it is rarely convenient or necessary to
consider all energies belonging to the system. Rather, we are far more interested in the change in internal energy, given some
transfer of work or heat. This can be expressed as:
ΔU = Q + Wmech + Wother . (13.1.3)
Q is heat added to a system and Wmech is the mechanical work performed by the surroundings due to pressure or volume changes in
the system. All other perturbations and energies added by other processes, such as an electric current introduced into an electronic
circuit, is summarized as the term Wextra.
We can calculate a small change in internal energy of the system by considering the infinitesimal amount of heat δQ added to the
system minus the infinitesimal amount of work δW done by the system:
dU = δQ − δW. (13.1.4)
The Sun and Internal Energy: Nuclear fusion in the sun converts nuclear potential energy into available internal energy and
keeps the temperature of the Sun very high.
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Key Points
Heat is a crucial concept that touches every aspect of our lives. James Clerk Maxwell set down important principles that couple
into the definition of heat.
The quantity of heat transfer can be directly measured or estimated indirectly through the science of calorimetry.
There are three modes of heat transfer: conduction, convection, and radiation.
If two objects at different temperature are brought together, energy will transfer from the hotter object to the cooler one until
both are at the same temperature. This transfer of energy is known as heat.
Heat should not be confused with temperature. Temperature describes the internal state of an object, while heat refers to the
energy transferred to or from the object.
Since heat is a form of energy, its SI unit is the joule. Other common units of heat energy include the calorie and kilocalorie,
equal to 4.186 and 4,186 joules, respectively.
Because heat and work both involve the transfer of energy, they can each produce the same effects. The concept of the
mechanical equivalent of heat was instrumental in establishing the principle of conservation of energy.
While a system does not contain ‘ heat,’ it does contain a total amount of energy called internal energy.
The internal energy is the energy necessary to create a system, minus the energy necessary to displace its surroundings.
Most of the time, we are interested in the change in internal energy rather than the total internal energy.
The first law of thermodynamics, dU = ΔQ − ΔW , describes small changes in internal energy.
Key Terms
heat transfer: The transmission of thermal energy via conduction, convection, or radiation.
calorimetry: The science of measuring the heat absorbed or evolved during the course of a chemical reaction or change of
state.
kilocalorie: A non-SI unit of energy equal to 1,000 calories or 4,186 joules; equal to the “calorie” or “Calorie” used in
nutritional labeling. Symbol: kcal.
thermal equilibrium: Two systems are in thermal equilibrium if they could transfer heat between each other, but don’t.
mechanical equivalent of heat: The work needed to produce the same effects as heat transfer.
internal energy: The sum of all energy present in the system, including kinetic and potential energy; equivalently, the energy
needed to create a system, excluding the energy necessary to displace its surroundings.
isolated system: A system that does not interact with its surroundings, that is, its total energy and mass stay constant.
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13.2: Specific Heat
learning objectives
Explain the enthalpy in a system with constant volume and pressure
Heat capacity (usually denoted by a capital C, often with subscripts), or thermal capacity, is the measurable physical quantity that
characterizes the amount of heat required to change a substance’s temperature by a given amount. In SI units, heat capacity is
expressed in units of joules per kelvin (J/K).
An object’s heat capacity (symbol C) is defined as the ratio of the amount of heat energy transferred to an object to the resulting
increase in temperature of the object.
Q
C = . (13.2.1)
ΔT
Heat capacity is an extensive property, so it scales with the size of the system. A sample containing twice the amount of substance
as another sample requires the transfer of twice as much heat (Q) to achieve the same change in temperature (ΔT). For example, if
it takes 1,000 J to heat a block of iron, it would take 2,000 J to heat a second block of iron with twice the mass as the first.
If the heat is added at constant volume, then the second term of this relation vanishes and one readily obtains
∂U ∂Q
( )V = ( )V = CV . (13.2.4)
∂T ∂T
This defines the heat capacity at constant volume, CV. Another useful quantity is the heat capacity at constant pressure, CP. With
the enthalpy of the system given by
H = U + PV, (13.2.5)
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∂H ∂Q
( )P = ( )P = CP . (13.2.7)
∂T ∂T
Specific Heat
The specific heat is an intensive property that describes how much heat must be added to a particular substance to raise its
temperature.
learning objectives
Summarize the quantitative relationship between heat transfer and temperature change
Specific Heat
The heat capacity is an extensive property that describes how much heat energy it takes to raise the temperature of a given system.
However, it would be pretty inconvenient to measure the heat capacity of every unit of matter. What we want is an intensive
property that depends only on the type and phase of a substance and can be applied to systems of arbitrary size. This quantity is
known as the specific heat capacity (or simply, the specific heat), which is the heat capacity per unit mass of a material.
Experiments show that the transferred heat depends on three factors: (1) The change in temperature, (2) the mass of the system, and
(3) the substance and phase of the substance. The last two factors are encapsulated in the value of the specific heat.
Heat Transfer and Specific Heat Capacity: The heat Q transferred to cause a temperature change depends on the magnitude of
the temperature change, the mass of the system, and the substance and phase involved. (a) The amount of heat transferred is
directly proportional to the temperature change. To double the temperature change of a mass m, you need to add twice the heat. (b)
The amount of heat transferred is also directly proportional to the mass. To cause an equivalent temperature change in a doubled
mass, you need to add twice the heat. (c) The amount of heat transferred depends on the substance and its phase. If it takes an
amount Q of heat to cause a temperature change ΔT in a given mass of copper, it will take 10.8 times that amount of heat to cause
the equivalent temperature change in the same mass of water assuming no phase change in either substance.
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Chemistry 10.2 Speci c Heat Capacity
Specific Heat Capacity: This lesson relates heat to a change in temperature. We discuss how the amount of heat needed for a
temperature change is dependent on mass and the substance involved, and that relationship is represented by the specific heat
capactiy of the substance, C.
The dependence on temperature change and mass are easily understood. Because the (average) kinetic energy of an atom or
molecule is proportional to the absolute temperature, the internal energy of a system is proportional to the absolute temperature and
the number of atoms or molecules. Since the transferred heat is equal to the change in the internal energy, the heat is proportional to
the mass of the substance and the temperature change. The transferred heat also depends on the substance so that, for example, the
heat necessary to raise the temperature is less for alcohol than for water. For the same substance, the transferred heat also depends
on the phase (gas, liquid, or solid).
The quantitative relationship between heat transfer and temperature change contains all three factors:
Q = mcΔT, (13.2.8)
where Q is the symbol for heat transfer, m is the mass of the substance, and ΔT is the change in temperature. The symbol c stands
for specific heat and depends on the material and phase.
The specific heat is the amount of heat necessary to change the temperature of 1.00 kg of mass by 1.00ºC. The specific heat c is a
property of the substance; its SI unit is J/(kg⋅K) or J/(kg⋅C). Recall that the temperature change (ΔT) is the same in units of kelvin
and degrees Celsius. Note that the total heat capacity C is simply the product of the specific heat capacity c and the mass of the
substance m, i.e.,
C C
C = mc or c = = , (13.2.9)
m ρV
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where ϱ is the density of the substance and V is its volume.
Values of specific heat must generally be looked up in tables, because there is no simple way to calculate them. Instead, they are
measured empirically. In general, the specific heat also depends on the temperature. The table below lists representative values of
specific heat for various substances. Except for gases, the temperature and volume dependence of the specific heat of most
substances is weak. The specific heat of water is five times that of glass and ten times that of iron, which means that it takes five
times as much heat to raise the temperature of water the same amount as for glass and ten times as much heat to raise the
temperature of water as for iron. In fact, water has one of the largest specific heats of any material, which is important for
sustaining life on Earth.
Specific Heats: Listed are the specific heats of various substances. These values are identical in units of cal/(g⋅C).3. cv at constant
volume and at 20.0ºC, except as noted, and at 1.00 atm average pressure. Values in parentheses are cp at a constant pressure of 1.00
atm.
Calorimetry
Calorimetry is the measurement of the heat of chemical reactions or physical changes.
learning objectives
Analyze the relationship between the gas constant for an ideal gas yield and volume
Calorimetry
Overview
Calorimetry is the science of measuring the heat of chemical reactions or physical changes. Calorimetry is performed with a
calorimeter. A simple calorimeter just consists of a thermometer attached to a metal container full of water suspended above a
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combustion chamber. The word calorimetry is derived from the Latin word calor, meaning heat. Scottish physician and scientist
Joseph Black, who was the first to recognize the distinction between heat and temperature, is said to be the founder of calorimetry.
Calorimetry requires that the material being heated have known thermal properties, i.e. specific heat capacities. The classical rule,
recognized by Clausius and by Kelvin, is that the pressure exerted by the calorimetric material is fully and rapidly determined
solely by its temperature and volume; this rule is for changes that do not involve phase change, such as melting of ice. There are
many materials that do not comply with this rule, and for them, more complex equations are required than those below.
Ice Calorimeter: The world’s first ice-calorimeter, used in the winter of 1782-83, by Antoine Lavoisier and Pierre-Simon Laplace,
to determine the heat evolved in variouschemical changes; calculations which were based on Joseph Black’s prior discovery of
latent heat. These experiments mark the foundation of thermochemistry.
where δQ is the increment of heat gained by the sample, CV is the heat capacity at constant volume, cv is the specific heat at
constant volume, and ΔT is the change in temperature.
Dividing the energy change by how many grams (or moles) of A were present gives its enthalpy change of reaction. This method is
used primarily in academic teaching as it describes the theory of calorimetry. It does not account for the heat loss through the
container or the heat capacity of the thermometer and container itself. In addition, the object placed inside the calorimeter shows
that the objects transferred their heat to the calorimeter and into the liquid, and the heat absorbed by the calorimeter and the liquid
is equal to the heat given off by the metals.
Constant-Pressure Calorimetry
A constant-pressure calorimeter measures the change in enthalpy of a reaction occurring in solution during which the atmospheric
pressure remains constant. An example is a coffee-cup calorimeter, which is constructed from two nested Styrofoam cups and a lid
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with two holes, allowing insertion of a thermometer and a stirring rod. The inner cup holds a known amount of a solute, usually
water, that absorbs the heat from the reaction. When the reaction occurs, the outer cup provides insulation. Then
WΔH
CP = (13.2.12)
MΔT
where Cp is the specific heat at constant pressure, ΔH is the enthalpy of the solution, ΔT is the change in temperature, W is the
mass of the solute, and M is the molecular mass of the solute. The measurement of heat using a simple calorimeter, like the coffee
cup calorimeter, is an example of constant-pressure calorimetry, since the pressure (atmospheric pressure) remains constant during
the process. Constant-pressure calorimetry is used in determining the changes in enthalpy occurring in solution. Under these
conditions the change in enthalpy equals the heat (Q=ΔH).
learning objectives
Explain how to derive the adiabatic index
This represents the dimensionless heat capacity at constant volume; it is generally a function of temperature due to intermolecular
forces. For moderate temperatures, the constant for a monoatomic gas is cv=3/2 while for a diatomic gas it is cv=5/2 (see ).
Macroscopic measurements on heat capacity provide information on the microscopic structure of the molecules.
Molecular internal vibrations: When a gas is heated, translational kientic energy of molecules in the gas will increase. In
addition, molecules in the gas may pick up many characteristic internal vibrations. Potential energy stored in these internal degrees
of freedom contributes to specific heat of the gas.
The heat capacity at constant pressure of 1 J·K−1 ideal gas is:
∂H
( )V = cp = cv + R (13.2.14)
∂T
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∂P ∂V
Cp − CV = T( )V , N( )p,N (13.2.15)
∂T ∂T
where the partial derivatives are taken at: constant volume and constant number of particles, and at constant pressure and constant
number of particles, respectively.
The heat capacity ratio or adiabatic index is the ratio of the heat capacity at constant pressure to heat capacity at constant volume. It
is sometimes also known as the isentropic expansion factor:
CP cp
γ = = (13.2.16)
CV cv
For an ideal gas, evaluating the partial derivatives above according to the equation of state, where R is the gas constant for an ideal
gas yields:
pV = RT (13.2.17)
∂P ∂V
Cp − CV = T( )V ( )p (13.2.18)
∂T ∂T
∂P ∂V
2
Cp − CV = −T( )V ( )p (13.2.19)
∂V ∂T
RT ∂P −RT −P
P = n→ ( )T = = (13.2.20)
2
V ∂V V V
2
RT ∂V R
2
V = n→ ( )p = (13.2.21)
2
P ∂T P
substituting:
2
∂P ∂V −P R
2
−T( )V ( )p = −T =R (13.2.22)
2
∂V ∂T V P
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Julius Robert Mayer: Julius Robert von Mayer (November 25, 1814 – March 20, 1878), a German physician and physicist, was
one of the founders of thermodynamics. He is best known for his 1841 enunciation of one of the original statements of the
conservation of energy (or what is now known as one of the first versions of the first law of thermodynamics): “Energy can be
neither created nor destroyed. ” In 1842, Mayer described the vital chemical process now referred to as oxidation as the primary
source of energy for any living creature. His achievements were overlooked and credit for the discovery of the mechanical
equivalent of heat was attributed to James Joule in the following year. von Mayer also proposed that plants convert light into
chemical energy.
CP − CV = R. (13.2.23)
It is a simple equation relating the heat capacities under constant temperature and under constant pressure.
learning objectives
Explain a bomb calorimeter is used to measure heat evolved in a combustion reaction
Calorimeters are designed to minimize energy exchange between the system being studied and its surroundings. They range from
simple coffee cup calorimeters used by introductory chemistry students to sophisticated bomb calorimeters used to determine the
energy content of food.
Calorimetry is used to measure amounts of heat transferred to or from a substance. To do so, the heat is exchanged with a calibrated
object (calorimeter). The change in temperature of the measuring part of the calorimeter is converted into the amount of heat (since
the previous calibration was used to establish its heat capacity ). The measurement of heat transfer using this approach requires the
definition of a system (the substance or substances undergoing the chemical or physical change) and its surroundings (the other
components of the measurement apparatus that serve to either provide heat to the system or absorb heat from the system).
Knowledge of the heat capacity of the surroundings, and careful measurements of the masses of the system and surroundings and
their temperatures before and after the process allows one to calculate the heat transferred as described in this section.
A calorimeter is a device used to measure the amount of heat involved in a chemical or physical process. For example, when an
exothermic reaction occurs in solution in a calorimeter, the heat produced by the reaction is absorbed by the solution, which
increases its temperature. When an endothermic reaction occurs, the heat required is absorbed from the thermal energy of the
solution, which decreases its temperature. The temperature change, along with the specific heat and mass of the solution, can then
be used to calculate the amount of heat involved in either case.
Coffee-Cup Calorimeters
General chemistry students often use simple calorimeters constructed from polystyrene cups. These easy-to-use “coffee cup”
calorimeters allow more heat exchange with their surroundings, and therefore produce less accurate energy values.
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A different type of calorimeter that operates at constant volume, colloquially known as a bomb calorimeter, is used to measure the
energy produced by reactions that yield large amounts of heat and gaseous products, such as combustion reactions. (The term
“bomb” comes from the observation that these reactions can be vigorous enough to resemble explosions that would damage other
calorimeters.) This type of calorimeter consists of a robust steel container (the “bomb”) that contains the reactants and is itself
submerged in water. The sample is placed in the bomb, which is then filled with oxygen at high pressure. A small electrical spark is
used to ignite the sample. The energy produced by the reaction is trapped in the steel bomb and the surrounding water. The
temperature increase is measured and, along with the known heat capacity of the calorimeter, is used to calculate the energy
produced by the reaction. Bomb calorimeters require calibration to determine the heat capacity of the calorimeter and ensure
accurate results. The calibration is accomplished using a reaction with a known q, such as a measured quantity of benzoic acid
ignited by a spark from a nickel fuse wire that is weighed before and after the reaction. The temperature change produced by the
known reaction is used to determine the heat capacity of the calorimeter. The calibration is generally performed each time before
the calorimeter is used to gather research data.
A 59.7 g piece of metal that had been submerged in boiling water was quickly transferred into 60.0 mL of water initially at
22.0 °C. The final temperature is 28.5 °C. Use these data to determine the specific heat of the metal. Use this result to identify
the metal.
Solution
Assuming perfect heat transfer, the heat given off by metal is the negative of the heat taken in by water, or:
cmetal × mmetal × (Tf ,metal − Ti,metal) = cwater × mwater × (Tf ,water − Ti,water) (13.2.25)
Noting that since the metal was submerged in boiling water, its initial temperature was 100.0 °C; and that for water, 60.0 mL =
60.0 g; we have:
o o o
(cmetal )(59.7g)(28.5 C − 100.0 C) = (4.18J/ g C)(60.0g)(28.5oC − 22.0oC) (13.2.26)
Solving this:
o o
−(4.184J/ g C)(60.0g)(6.5 C)
o
cmetal = = 0.38J/ g C (13.2.27)
o
(59.7g)(−71.5 C)
Our experimental specific heat is closest to the value for copper (0.39 J/g °C), so we identify the metal as copper.
Key Points
Heat capacity is the measurable physical quantity that characterizes the amount of heat required to change a substance’s
temperature by a given amount. It is measured in joules per Kelvin and given by.
The heat capacity is an extensive property, scaling with the size of the system.
The heat capacity of most systems is not constant (though it can often be treated as such). It depends on the temperature,
pressure, and volume of the system under consideration.
Unlike the total heat capacity, the specific heat capacity is independent of mass or volume. It describes how much heat must be
added to a unit of mass of a given substance to raise its temperature by one degree Celsius. The units of specific heat capacity
are J/(kg °C) or equivalently J/(kg K).
The heat capacity and the specific heat are related by C = cm or c = . C
The mass m, specific heat c, change in temperature ΔT, and heat added (or subtracted) Q are related by the equation:
Q = mcΔT .
Values of specific heat are dependent on the properties and phase of a given substance. Since they cannot be calculated easily,
they are empirically measured and available for reference in tables.
A calorimeter is used to measure the heat generated (or absorbed) by a physical change or chemical reaction. The science of
measuring these changes is known as calorimetry.
In order to do calorimetry, it is crucial to know the specific heats of the substances being measured.
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Calorimetry can be performed under constant volume or constant pressure. The type of calculation done depends on the
conditions of the experiment.
The specific heat at constant volume for a gas is given as ( ) = c .
∂U
∂T
V v
∂T
V p v
The heat capacity ratio (or adiabatic index ) is the ratio of the heat capacity at constant pressure to heat capacity at constant
volume.
Calorimetry is used to measure amounts of heat transferred to or from a substance.
A calorimeter is a device used to measure the amount of heat involved in a chemical or physical process.
This means that the amount of heat produced or consumed in the reaction equals the amount of heat absorbed or lost by the
solution.
Key Terms
heat capacity: The amount of heat energy needed to raise the temperature of an object or unit of matter by one degree Celsius;
in units of joules per kelvin (J/K).
enthalpy: the total amount of energy in a system, including both the internal energy and the energy needed to displace its
environment
specific heat capacity: The amount of heat that must be added (or removed) from a unit mass of a substance to change its
temperature by one degree Celsius. It is an intensive property.
constant-pressure calorimeter: An instrument used to measure the heat generated during changes that do not involve changes
in pressure.
calorimeter: An apparatus for measuring the heat generated or absorbed by either a chemical reaction, change of phase or some
other physical change.
constant-volume calorimeter: An instrument used to measure the heat generated during changes that do not involve changes in
volume.
Fundamental Thermodynamic Relation: In thermodynamics, the fundamental thermodynamic relation expresses an
infinitesimal change in internal energy in terms of infinitesimal changes in entropy, and volume for a closed system in thermal
equilibrium in the following way: dU=TdS-PdV. Here, U is internal energy, T is absolute temperature, S is entropy, P is pressure
and V is volume.
adiabatic index: The ratio of the heat capacity at constant pressure to heat capacity at constant volume.
specific heat: The ratio of the amount of heat needed to raise the temperature of a unit mass of substance by a unit degree to the
amount of heat needed to raise that of the same mass of water by the same amount.
heat of reaction: The enthalpy change in a chemical reaction; the amount of heat that a systems gives up to its surroundings so
it can return to its initial temperature.
combustion: A process where two chemicals are combined to produce heat.
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13.3: Phase Change and Latent Heat
learning objectives
Describe the latent heat as a form of energy
Latent Heat
Previously, we have discussed temperature change due to heat transfer. No temperature change occurs from heat transfer if ice
melts and becomes liquid water (i.e., during a phase change). For example, consider water dripping from icicles melting on a roof
warmed by the Sun. Conversely, water freezes in an ice tray cooled by lower-temperature surroundings.
Melting Icicle: Heat from the air transfers to the ice causing it to melt.
Energy is required to melt a solid because the cohesive bonds between the molecules in the solid must be broken apart so that the
molecules can move around at comparable kinetic energies; thus, there is no rise in temperature. Similarly, energy is needed to
vaporize a liquid, because molecules in a liquid interact with each other via attractive forces. There is no temperature change until a
phase change is complete. The temperature of a glass of lemonade initially at 0 ºC stays at 0 ºC until all the ice has melted.
Conversely, energy is released during freezing and condensation, usually in the form of thermal energy. Work is done by cohesive
forces when molecules are brought together. The corresponding energy must be given off (dissipated) to allow them to stay
together.
The energy involved in a phase change depends on two major factors: the number and strength of bonds or force pairs. The number
of bonds is proportional to the number of molecules and thus to the mass of the sample. The strength of forces depends on the type
of molecules. The heat Q required to change the phase of a sample of mass m is given by
Q = mLf (melting or freezing)
Q = mLv (evaporating or condensing)
where the latent heat of fusion, Lf, and latent heat of vaporization, Lv, are material constants that are determined experimentally.
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Phase Transitions: (a) Energy is required to partially overcome the attractive forces between molecules in a solid to form a liquid.
That same energy must be removed for freezing to take place. (b) Molecules are separated by large distances when going from
liquid to vapor, requiring significant energy to overcome molecular attraction. The same energy must be removed for condensation
to take place. There is no temperature change until a phase change is complete.
Latent heat is an intensive property measured in units of J/kg. Both Lf and Lv depend on the substance, particularly on the strength
of its molecular forces as noted earlier. Lf and Lv are collectively called latent heat coefficients. They are latent, or hidden, because
in phase changes, energy enters or leaves a system without causing a temperature change in the system; so, in effect, the energy is
hidden. Note that melting and vaporization are endothermic processes in that they absorb or require energy, while freezing and
condensation are exothermic process as they release energy.
13.3.2 https://phys.libretexts.org/@go/page/14517
Heating Up Ice: Andrew Vanden Heuvel explores latent heat while trying to cool down his soda.
Significant amounts of energy are involved in phase changes. Let us look, for example, at how much energy is needed to melt a
kilogram of ice at 0º C to produce a kilogram of water at 0°C. Using the equation for a change in temperature and the value for
water (334 kJ/kg), we find that Q=mLf=(1.0kg)(334kJ/kg)=334kJ is the energy to melt a kilogram of ice. This is a lot of energy as
it represents the same amount of energy needed to raise the temperature of 1 kg of liquid water from 0ºC to 79.8ºC. Even more
energy is required to vaporize water; it would take 2256 kJ to change 1 kg of liquid water at the normal boiling point (100ºC at
atmospheric pressure) to steam (water vapor). This example shows that the energy for a phase change is enormous compared to
energy associated with temperature changes without a phase change.
Phase changes can have an enormous stabilizing effect (see figure below). Consider adding heat at a constant rate to a sample of ice
initially at -20 ºC. Initially the temperature of the ice rises linearly, absorbing heat at a constant rate of 0.50 cal/g⋅C until it reaches
0 ºC. Once at this temperature, the ice begins to melt until all the entire sample has melted, absorbing a total of 79.8 cal/g of heat.
The temperature remains constant at 0 ºC during this phase change. Once all the ice has melted, the temperature of the liquid water
rises, absorbing heat at a new constant rate of 1.00 cal/g⋅C (remember that specific heats are dependent on phase). At 100ºC, the
water begins to boil and the temperature again remains constant until the water absorbs 539 cal/g of heat to complete this phase
change. When all the liquid has become steam, the temperature rises again, absorbing heat at a rate of 0.482 cal/g⋅C.
Heating and Phase Changes of Water: A graph of temperature versus energy added. The system is constructed so that no vapor
evaporates while ice warms to become liquid water, and so that, when vaporization occurs, the vapor remains in of the system. The
long stretches of constant temperature values at 0ºC and 100ºC reflect the large latent heat of melting and vaporization,
respectively.
A phase change we have neglected to mention thus far is sublimation, the transition of solid directly into vapor. The opposite case,
where vapor transitions directly into a solid, is called deposition. Sublimation has its own latent heat Ls and can be used in the
same way as Lvand Lf.
Key Points
Energy is required to change the phase of a substance, such as the energy to break the bonds between molecules in a block of
ice so it may melt.
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During a phase change energy my be added or subtracted from a system, but the temperature will not change. The temperature
will change only when the phase change has completed.
The heat Q required to change the phase of a sample of mass m is given by Q = mL (melting or freezing) and Q = mL
f v
(evaporating or condensing), where L and L are the latent heat of fusion and the latent heat of vaporization, respectively.
f v
Key Terms
latent heat of fusion: the energy required to transition one unit of a substance from solid to liquid; equivalently, the energy
liberated when one unit of a substance transitions from liquid to solid.
latent heat of vaporization: the energy required to transition one unit of a substance from liquid to vapor; equivalently, the
energy liberated when one unit of a substance transitions from vapor to liquid.
sublimation: the transition of a substance from the solid phase directly to the vapor state such that it does not pass through the
intermediate, liquid phase
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13.4: Methods of Heat Transfer
learning objectives
Assess why particular characteristics are necessary for effective conduction
Conduction
Conduction is the transfer of heat through stationary matter by physical contact. (The matter is stationary on a macroscopic scale—
we know there is thermal motion of the atoms and molecules at any temperature above absolute zero.) Heat transferred from an
electric stove to the bottom of a pot is an example of conduction.
Some materials conduct thermal energy faster than others. For example, the pillow in your room may the same temperature as the
metal doorknob, but the doorknob feels cooler to the touch. In general, good conductors of electricity (metals like copper,
aluminum, gold, and silver) are also good heat conductors, whereas insulators of electricity (wood, plastic, and rubber) are poor
heat conductors.
Microscopic Description of Conduction
On a microscopic scale, conduction occurs as rapidly moving or vibrating atoms and molecules interact with neighboring particles,
transferring some of their kinetic energy. Heat is transferred by conduction when adjacent atoms vibrate against one another, or as
electrons move from one atom to another. Conduction is the most significant means of heat transfer within a solid or between solid
objects in thermal contact. Conduction is greater in solids because the network of relatively close fixed spatial relationships
between atoms helps to transfer energy between them by vibration.
Fluids and gases are less conductive than solids. This is due to the large distance between atoms in a fluid or (especially) a gas:
fewer collisions between atoms means less conduction.
Microscopic Illustration of Conduction: The molecules in two bodies at different temperatures have different average kinetic
energies. Collisions occurring at the contact surface tend to transfer energy from high-temperature regions to low-temperature
regions. In this illustration, a molecule in the lower temperature region (right side) has low energy before collision, but its energy
increases after colliding with the contact surface. In contrast, a molecule in the higher temperature region (left side) has high energy
before collision, but its energy decreases after colliding with the contact surface.
The (average) kinetic energy of a molecule in the hot body is higher than in the colder body. If two molecules collide, an energy
transfer from the hot to the cold molecule occurs (see the above figure). The cumulative effect from all collisions results in a net
flux of heat from the hot body to the colder body. The heat flux thus depends on the temperature difference T = T − T hot .
cold
Therefore, you will get a more severe burn from boiling water than from hot tap water. Conversely, if the temperatures are the
same, the net heat transfer rate falls to zero, and equilibrium is achieved. Owing to the fact that the number of collisions increases
with increasing area, heat conduction depends on the cross-sectional area. If you touch a cold wall with your palm, your hand cools
faster than if you just touch it with your fingertip.
Factors Affecting the Rate of Heat Transfer Through Conduction
In addition to temperature and cross-sectional area, another factor affecting conduction is the thickness of the material through
which the heat transfers. Heat transfer from the left side to the right side is accomplished by a series of molecular collisions. The
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thicker the material, the more time it takes to transfer the same amount of heat. If you get cold during the night, you may retrieve a
thicker blanket to keep warm.
Effect of Thickness on Heat Conduction: Heat conduction occurs through any material, represented here by a rectangular bar.
The temperature of the material is T on the left and T on the right, where T is greater than T . The rate of heat transfer by
2 1 2 1
conduction is directly proportional to the surface area A , the temperature difference T − T , and the substance’s conductivity kk.
2 1
Joules per second (Watts), k is the thermal conductivity of the material, A and d are its surface area and thickness, and (T − T )
2 1
Convection
Convection is the heat transfer by the macroscopic movement of a fluid, such as a car’s engine kept cool by the water in the cooling
system.
learning objectives
Illustrate the mechanisms of convection with phase change
Example 13.4.1:
Calculating Heat Transfer by Convection: Convection of Air Through the Walls of a House.
Most houses are not airtight: air goes in and out around doors and windows, through cracks and crevices, following wiring to
switches and outlets, and so on. The air in a typical house is completely replaced in less than an hour.
Suppose that a moderately-sized house has inside dimensions 12.0 m × 18.0 m × 3.00 m high, and that all air is replaced in
30.0 min. Calculate the heat transfer per unit time in watts needed to warm the incoming cold air by 10.0 ºC, thus replacing the
heat transferred by convection alone.
Strategy:
Q
Heat is used to raise the temperature of air so that Q = mcΔT . The rate of heat transfer is then , where t is the time for air
t
turnover. We are given that ΔT is 10.0ºC, but we must still find values for the mass of air and its specific heat before we can
calculate QQ. The specific heat of air is a weighted average of the specific heats of nitrogen and oxygen, which is
c = cp ≅1000 J/kg ⋅ C (note that the specific heat at constant pressure must be used for this process).
Solution
(1) Determine the mass of air from its density and the given volume of the house. The density is given from the density ρ and
the volume m = ρV = (1.29 kg/m3)(12.0 m × 18.0 m × 3.00 m) = 836 kg
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(2) Calculate the heat transferred from the change in air temperature: Q = mcΔT so that
6
Q = (836 kg)(1000 J/kg ⋅ ∘C)(10 ∘ C) = 8.36 × 10 J
(3) Calculate the heat transfer from the heat Q and the turnover time t . Since air is turned over in t = 0.500 h = 1800 s , the
6
Q
heat transferred per unit time is =
t
8.36×10
1800 s
= 4.64 kW .
J
This rate of heat transfer is equal to the power consumed by about forty-six 100-W light bulbs.
Newly constructed homes are designed for a turnover time of 2 hours or more, rather than 30 minutes for the house of this
example. Weather stripping, caulking, and improved window seals are commonly employed. More extreme measures are
sometimes taken in very cold (or hot) climates to achieve a tight standard of more than 6 hours for one air turnover. Still longer
turnover times are unhealthy, because a minimum amount of fresh air is necessary to supply oxygen for breathing and to dilute
household pollutants. The term used for the process by which outside air leaks into the house from cracks around windows,
doors, and the foundation is called “air infiltration.”
Convection
Convection (illustrated in ) is the concerted, collective movement of ensembles of molecules within fluids (e.g., liquids, gases).
Convection of mass cannot take place in solids, since neither bulk current flows nor significant diffusion can occur in solids.
Instead heat diffusion in solids is called heat conduction, which we’ve just reviewed.
Convection in a Pot of Water: Convection plays an important role in heat transfer inside this pot of water. Once conducted to the
inside, heat transfer to other parts of the pot is mostly by convection. The hotter water expands, decreases in density, and rises to
transfer heat to other regions of the water, while colder water sinks to the bottom. This process keeps repeating.
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Convection and Insulation
Although air can transfer heat rapidly by convection, it is a poor conductor and thus a good insulator. The amount of available
space for airflow determines whether air acts as an insulator or conductor. The space between the inside and outside walls of a
house, for example, is about 9 cm (3.5 in)—large enough for convection to work effectively. The addition of wall insulation
prevents airflow, so heat loss (or gain) is decreased. Similarly, the gap between the two panes of a double-paned window is about 1
cm, which prevents convection and takes advantage of air’s low conductivity to prevent greater loss. Fur, fiber and fiberglass also
take advantage of the low conductivity of air by trapping it in spaces too small to support convection. In animals, fur and feathers
are lightweight and thus ideal for their protection.
Convection and Phase Changes
Some interesting phenomena happen when convection is accompanied by a phase change. It allows us to cool off by sweating, even
if the temperature of the surrounding air exceeds body temperature. Heat from the skin is required in order for sweat to evaporate
from the skin, but without air flow the air becomes saturated and evaporation stops. Air flow caused by convection replaces the
saturated air by dry air and thus evaporation continues.
Another important example of the combination of phase change and convection occurs when water evaporates from the ocean.
Heat is removed from the ocean when water evaporates. If the water vapor condenses in liquid droplets as clouds form, heat is
released in the atmosphere (this heat release is latent heat) . Thus, an overall transfer of heat from the ocean to the atmosphere
occurs. This process is the driving power behind thunderheads—great cumulus clouds that rise as much as 20.0 km into the
stratosphere. Water vapor carried in by convection condenses, releasing tremendous amounts of energy, and this energy allows air
to become more buoyant (warmer than its surroundings) and rise. As the air continues to rise, more condensation occurs, which in
turn drives the cloud even higher. Such a mechanism is called positive feedback, since the process reinforces and accelerates itself.
These systems sometimes produce violent storms with lightning and hail, and constitute the mechanism that drives hurricanes.
Cumulus Clouds: Cumulus clouds are caused by water vapor that rises because of convection. The rise of clouds is driven by a
positive feedback mechanism.
Radiation
Radiation is the transfer of heat through electromagnetic energy
learning objectives
Explain how the energy of electromagnetic radiation corresponds with wavelength
Radiation
You can feel heat transfer from a fire or the Sun. Yet the space between Earth and the Sun is largely empty, without any possibility
of heat transfer by convection or conduction. Similarly, you can tell that an oven is hot without touching it or looking inside—it just
warms you as you walk by.
In these examples, heat is transferred by radiation. The hot body emits electromagnetic waves that are absorbed by our skin, and no
medium is required for them to propagate. We use different names for electromagnetic waves of different wavelengths: radio
waves, microwaves, infrared radiation, visible light, ultraviolet radiation, X-rays, and gamma rays .
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Radiation from a Fire: Most of the heat transfer from this fire to the observers is through infrared radiation. The visible light,
although dramatic, transfers relatively little thermal energy. Convection transfers energy away from the observers as hot air rises,
while conduction is negligibly slow here. Skin is very sensitive to infrared radiation so that you can sense the presence of a fire
without looking at it directly.
The energy of electromagnetic radiation depends on its wavelength (color) and varies over a wide range; a smaller wavelength (or
higher frequency) corresponds to a higher energy. We can write this as:
hc
E = hf = (13.4.1)
λ
Radiation Spectrum: (a) A graph of the spectra of electromagnetic waves emitted from an ideal radiator at three different
temperatures. The intensity or rate of radiation emission increases dramatically with temperature, and the spectrum shifts toward
the visible and ultraviolet parts of the spectrum. The shaded portion denotes the visible part of the spectrum. It is apparent that the
shift toward the ultraviolet with temperature makes the visible appearance shift from red to white to blue as temperature increases.
(b) Note the variations in color corresponding to variations in flame temperature.
Heat Transfer
All objects absorb and emit electromagnetic radiation. The rate of heat transfer by radiation is largely determined by the color of
the object. Black is the most effective, and white the least. People living in hot climates generally avoid wearing black clothing, for
instance. Similarly, black asphalt in a parking lot will be hotter than the adjacent gray sidewalk on a summer day, because black
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absorbs better than gray. The reverse is also true—black radiates better than gray. Thus, on a clear summer night the asphalt will be
colder than the gray sidewalk because black radiates energy more rapidly than gray.
An ideal radiator, often called a blackbody, is the same color as an ideal absorber, and captures all the radiation that falls on it. In
contrast, white is a poor absorber and also a poor radiator. A white object reflects all radiation, like a mirror. (A perfect, polished
white surface is mirror-like in appearance, and a crushed mirror looks white. )
There is a clever relation between the temperature of an ideal radiator and the wavelength at which it emits the most radiation. It is
called Wien’s displacement law and is given by:
λm axT = b (13.4.2)
Good and Poor Radiators: A black object is a good absorber and a good radiator, while a white (or silver) object is a poor
absorber and a poor radiator. It is as if radiation from the inside is reflected back into the silver object, whereas radiation from the
inside of the black object is “absorbed” when it hits the surface and finds itself on the outside and is strongly emitted.
The rate of heat transfer by emitted radiation is determined by the Stefan-Boltzmann law of radiation:
Q
4
= σeAT (13.4.3)
t
where σ = 5.67 × 10
−8 J
4
is the Stefan-Boltzmann constant, A is the surface area of the object, and T is its absolute
s⋅ m2 ⋅ K
temperature in kelvin. The symbol e stands for the emissivity of the object, which is a measure of how well it radiates. An ideal jet-
black (or blackbody) radiator has e=1e=1, whereas a perfect reflector has e = 0 . Real objects fall between these two values. For
example, tungsten light bulb filaments have an ee of about 0.5, and carbon black (a material used in printer toner), has the (greatest
known) emissivity of about 0.99.
The radiation rate is directly proportional to the fourth power of the absolute temperature—a remarkably strong temperature
dependence. Furthermore, the radiated heat is proportional to the surface area of the object. If you knock apart the coals of a fire,
there is a noticeable increase in radiation due to an increase in radiating surface area.
Net Rate of Heat Transfer
The net rate of heat transfer by radiation (absorption minus emission) is related to both the temperature of the object and that of its
surroundings. Assuming that an object with a temperature T is surrounded by an environment with uniform temperature T , the
1 2
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where e is the emissivity of the object alone. In other words, it does not matter whether the surroundings are white, gray, or black;
the balance of radiation into and out of the object depends on how well it emits and absorbs radiation. When T > T , the quantity
2 1
Q
net
t
is positive; that is, the net heat transfer is from hotter objects to colder objects.
Key Points
On a microscopic scale, conduction occurs as rapidly moving or vibrating atoms and molecules interact with neighboring
particles, transferring some of their kinetic energy.
Conduction is the most significant form of heat transfer within a solid object or between solids in thermal contact.
Conduction is most significant in solids, and less though in liquids and gases, due to the space between molecules.
The rate of heat transfer by conduction is dependent on the temperature difference, the size of the area in contact, the thickness
of the material, and the thermal properties of the material(s) in contact.
Convection is driven by the large scale flow of matter in fluids. Solids cannot transport heat through convection.
Natural convection is driven by buoyant forces: hot air rises because density decreases as temperature increases. This principle
applies equally with any fluid.
Convection can transport heat much more efficiently than conduction. Air is a poor conductor and a good insulator if the space
is small enough to prevent convection.
Convection often accompanies phase changes, such as when sweat evaporates from your body. This mass flow during
convection allows humans to cool off even if the surrounding air’s temperature exceeds the body temperature.
The energy of electromagnetic radiation depends on the wavelength (color) and varies over a wide range: a smaller wavelength
(or higher frequency) corresponds to a higher energy.
All objects emit and absorb electromagnetic energy. The color of an object is related emissivity, or its efficiency of radiating
away energy. Black is the most effective while white is the least effective (e = 1 and e = 0 , respectively).
An ideal radiator, often called a blackbody, is the same color as an ideal absorber and captures all the radiation that falls on it.
Q
The rate of heat transfer by emitted radiation is determined by the Stefan-Boltzmann law of radiation: = σAT where
t
4
σ = 5.67 × 10
−8 J
is the Stefan-Boltzmann constant, A is the surface area of the object, and T is its absolute
4
s⋅ m2 ⋅ K
temperature in kelvin.
The net rate of heat transfer is related to the temperature of the object and the temperature of its surroundings. The larger the
difference, the higher the net heat flux.
The temperature of an object is very significant, because the radiation emitted is proportional to this quantity to the fourth
power.
Key Terms
thermal conductivity: the measure of a material’s ability to conduct heat
natural convection: A method for heat transport. A fluid surrounding a heat source receives heat, becomes less dense and rises.
The surrounding, cooler fluid then moves to replace it. This cooler fluid is then heated and the process continues, forming a
convection current.
positive feedback: a feedback loop in which the output of a system is amplified with a net positive gain each cycle.
blackbody: A theoretical body, approximated by a hole in a hollow black sphere, that absorbs all incident electromagnetic
radiation and reflects none; it has a characteristic emission spectrum.
emissivity: The energy-emitting propensity of a surface, usually measured at a specific wavelength.
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13.5: Global Warming
learning objectives
Describe effect of greenhouse gases on the Earth’s equilibrium temperature
Radiation is a natural process of heat transfer; everything is constantly radiating heat. When a system is at equilibrium, we do not
notice this exchange of energy because every object absorbs and emits precisely the amount of energy it receives, but this
equilibrium is dynamic–the exchange is always taking place. In general, as the temperature of an object increases, it emits energy
more rapidly by emitting photons at a faster rate, and by emitting photons of a greater average energy.
If we approximate the Earth as a perfect absorber and emitter of the radiation received from the sun (called a blackbody), we would
expect the Earth to be at an average temperature of 5°C, rather than the 14°C which we observe. The extra temperature comes from
the Earth and it’s atmosphere selectively absorbing certain wavelengths of radiation while reflecting other wavelengths. The 9°C
discrepancy is due to the greenhouse effect. The gases of the atmosphere are “selective absorbers”; energy in the visible part of the
electromagnetic spectrum passes through the atmosphere directly to the Earth’s surface (with some reflection occurring as well).
The Earth absorbs this energy and then re-emits radiation in the infrared portion of the spectrum. The gases in the atmosphere,
primarily CO2 and water vapor are highly absorbent in the infrared part of the spectrum. The atmosphere absorbs the infrared
radiation from the Earth, preventing it from escaping to space. Because all objects are continually emitting radiation, the
atmosphere (having absorbed the Earth’s radiation) then emits radiation, some of which is then reabsorbed by the Earth’s surface.
Thus the greenhouse effect is a continuous cycle of absorption and emission of energy between the Earth and atmosphere. This
causes the Earth’s atmosphere and surface to be warmer than otherwise expected.
Radiative Transfer
The greenhouse effect is a phenomenon of radiative transfer, the process by which the energy of light waves is exchanged in matter.
Radiative transfer dictates what energy is reflected, absorbed, and emitted.
The greenhouse effect: A summary of the heat transfer in the Earth’s atmosphere.
Atmospheric Absorbers
The radiative transfer properties of atmospheric chemicals depend on the energy of the radiation (both for absorption and
emission), and those properties are unique to each chemical. In the context of global warming, we find that it is important to
consider both how chemicals and particles reflect sunlight and how they absorb energy radiated from the Earth. Atmospheric
reflecters, notably sulfates and nitrates, reflect and scatter light before it ever hits the surface of the Earth, effectively reducing the
power that the Earth receives. On the other hand, greenhouse gases such as carbon dioxide (CO2), methane (CH4), and nitrous
oxide (N2O) are characteristically strong absorbers of the energy radiated by the Earth’s surface. They absorb so strongly because
they typically exhibit resonant absorption behavior in the same energy range as the radiation emitted by the earth. These trap heat
before it leaves the Earth, insulating the Earth and increasing the Earth’s equilibrium temperature.
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radiates the same energy that it receives. (If greenhouse gases increase, then temperature increases, and higher temperatures cause
the Earth to radiate more, compensating for the greater energy absorbed in the atmosphere. ) With light radiation, we find the
entropy carried by the photons decreases as temperature increases, S∝1TS∝1T. Since the Earth is cooler than the Sun, as the Earth
absorbs radiation from the Sun and re-emits radiation from the Earth’s surface, there is a net production of entropy. This net
production of entropy allows energy to be stored and reused, particularly in the form of chemical energy stored by cells through
photosynthesis.
Key Points
The atmosphere can both absorb and reflect radiation, either increasing or decreasing the Earth’s average temperature.
All bodies naturally emit radiation; warmer bodies emit more photons, with higher average energies.
In absorbing sunlight and emitting infrared radiation, the Earth produces entropy.
Key Terms
selective absorber: An object that will absorb radiation over a particular set of wavelengths but will not (is transparent) at other
wavelengths.
greenhouse effect: The process by which a planet is warmed by its atmosphere.
greenhouse gas: Any gas, such as carbon dioxide or CFCs, that contributes to the greenhouse effect when released into the
atmosphere.
radiative transfer: The transfer of radiation (energy) leaving one object and being absorbed by another.
blackbody: An object that is a perfect absorber and emitter of radiation.
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13.6: Phase Equilbrium
learning objectives
Illustrate the causes for evaporation near the surface of a liquid
Evaporation is a type of vaporization of a liquid that only occurs on the liquid’s surface. Usually, the molecules in a glass of water
do not have enough heat energy to escape from the liquid. With sufficient heat, however, the liquid would quickly turn into vapor.
When the molecules collide, they transfer energy to each other in varying degrees. Sometimes the transfer is so one-sided for a
molecule near the surface that it achieves enough energy to escape the liquid.
Three key parts to evaporation are heat, atmospheric pressure (determines the percent humidity) and air movement. For molecules
of a liquid to evaporate, they must be located near the surface, be moving in the proper direction, and have sufficient kinetic energy
to overcome liquid- phase intermolecular forces. When only a small proportion of the molecules meet these criteria, the rate of
evaporation is low. Since the kinetic energy of a molecule is proportional to its temperature, evaporation proceeds more quickly at
higher temperatures. As the faster-moving molecules escape, the remaining molecules have lower average kinetic energy, and the
temperature of the liquid decreases. This phenomenon is also called evaporative cooling. This is why evaporating sweat cools the
human body. Evaporation also tends to proceed more quickly with higher flow rates between the gaseous and liquid phases and in
liquids with higher vapor pressure. For example, laundry on a clothes line will dry (by evaporation) more rapidly on a windy day
than on a still day.
Vapor Pressure Diagram: (a) Because of the distribution of speeds and kinetic energies, some water molecules can break away to
the vapor phase even at temperatures below the ordinary boiling point. (b) If the container is sealed, evaporation will continue until
there is enough vapor density for the condensation rate to equal the evaporation rate. This vapor density and the partial pressure it
creates are the saturation values. They increase with temperature and are independent of the presence of other gases, such as air.
They depend only on the vapor pressure of water.
Evaporation is an essential part of the water cycle. The sun (solar energy) drives evaporation of water from oceans, lakes, moisture
in the soil, and other sources of water. In hydrology, evaporation and transpiration (which involves evaporation within plant
stomata) are collectively termed evapotranspiration. Evaporation of water occurs when the surface of the liquid is exposed,
allowing molecules to escape and form water vapor; this vapor can then rise and form clouds.
learning objectives
Explain how a substance can have multiple distinct phases in the same environment
Phase Equilibrium
Left to equilibration, many compositions will form a uniform single phase, but depending on the temperature and pressure even a
single substance may separate into two or more distinct phases. Within each phase, the properties are uniform but between the two
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phases properties differ.
Water in a closed jar with an air space over it forms a two phase system. Most of the water is in the liquid phase, where it is held by
the mutual attraction of water molecules. Even at equilibrium, molecules are constantly in motion and, once in a while, a molecule
in the liquid phase gains enough kinetic energy to break away from the liquid phase and enter the gas phase. Likewise, every once
in a while a vapor molecule collides with the liquid surface and condenses into the liquid. At equilibrium, evaporation and
condensation processes exactly balance and there is no net change in the volume of either phase.
At room temperature and pressure, the water jar reaches equilibrium when the air over the water has a humidity of about 3%. This
percentage increases as the temperature goes up. At 100 °C and atmospheric pressure, equilibrium is not reached until the air is
100% water. If the liquid is heated a little over 100 °C, the transition from liquid to gas will occur not only at the surface, but
throughout the liquid volume: the water boils.
The Earth’s atmosphere is not unchanging. The water vapor in it changes phases. It is in a phase equilibrium. Collisions between
water molecules in the atmosphere allows some to condense and some to remain in vapor. Similarly, several lighter gases can
escape the gravitational field entirely.
Key Points
Evaporation turns liquids into gas.
Evaporation can take place at temperatures below boiling point since the molecules in the liquid have different energies.
As the molecules in a liquid collide, some achieve higher energies, allowing them to escape. This process lowers the energy of
the remaining molecules and is the source of cooling in evaporating liquids.
The atmosphere is made of gases in a phase equilibrium.
As molecules in the atmosphere collide, they gain and lose energy.
As water evaporates from the surface of the earth, water condenses in the atmosphere.
Key Terms
Vaporization: a conversion of a solid or a liquid into a gas
Evaporation: The process of a liquid converting to the gaseous state.
condensation: The conversion of a gas to a liquid; the condensate so formed
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
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CHAPTER OVERVIEW
14: Thermodynamics
Topic hierarchy
14.1: Introduction
14.2: The First Law of Thermodynamics
14.3: The Second Law of Thermodynamics
14.4: Entropy
14.5: The Third Law of Thermodynamics
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1
14.1: Introduction
learning objectives
Analyze the necessity to exclude energy transferred between system as heat from mechanical work
Work
In thermodynamics, work performed by a closed system is the energy transferred to another system that is measured by mechanical
constraints on the system. Thermodynamic work encompasses mechanical work (gas expansion, ) plus many other types of work,
such as electrical. As such, thermodynamic work is a generalization of the concept of mechanical work in mechanics. It necessarily
excludes energy transferred between systems as heat, which is modeled distinctly in thermodynamics. For closed systems, energy
changes in a system other than as work transfer are as heat.
Fig 1: An isobaric expansion of a gas requires heat transfer during the expansion to keep the pressure constant. Since pressure is
constant, the work done is PΔV.
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A Review of the Zeroth Law
Zeroth law justifies the use of thermodynamic temperature, defined as the shared temperature of three designated systems at
equilibrium.
learning objectives
Discuss how the Zeroth Law of Thermodynamics justifies the use of thermodynamic temperature
Thermodynamics and PV Diagrams: A brief introduction to the zeroth and 1st laws of thermodynamics as well as PV diagrams
for students.
14.1.2 https://phys.libretexts.org/@go/page/14522
Thermometer: A thermometer calibrated in degrees Celsius
The Zeroth Law of Thermodynamics states: If two systems, A and B, are in thermal equilibrium with each other, and B is in
thermal equilibrium with a third system, C, then A is also in thermal equilibrium with C.
This law was postulated in the 1930s, after the first and second laws of thermodynamics had been developed and named. It is called
the “zeroth” law because it comes logically before the first and second laws (discussed in Atoms on the 1st and 2nd laws).
Two systems are in thermal equilibrium if they could transfer heat between each other, but don’t. Indeed, experiments have shown
that if two systems, A and B, are in thermal equilibrium with each other, and B is in thermal equilibrium with a third system C, then
A is also in thermal equilibrium with C. This conclusion may seem obvious, because all three have the same temperature, but
zeroth law is basic to thermodynamics. Zeroth law justifies the use of thermodynamic temperature: the common “label” that the
three systems in the definition above share is defined as the temperature of the systems.
Temperature
Thermometers actually take their own temperature, not the temperature of the object they are measuring. This raises the question of
how we can be certain that a thermometer measures the temperature of the object with which it is in contact. The answer lies in the
fact that any two systems placed in thermal contact (meaning heat transfer can occur between them) will reach the same
temperature. That is, heat will flow from the hotter object to the cooler one until they reach exactly the same temperature. The
objects are then in thermal equilibrium, and no further changes will occur. The systems interact and change because their
temperatures differ, and the changes stop once their temperatures are the same. Thus, if enough time is allowed for this transfer of
heat to run its course, the temperature a thermometer registers does represent the system with which it achieves thermal
equilibrium.
Key Points
Thermodynamic work is a generalization of the concept of mechanical work in mechanics.
For closed systems, energy changes in a system other than as work transfer are as heat.
Work in thermodynamics is a quite organized process (as in gas expansion), involving a macroscopic force exerted through a
distance.
The zeroth law of thermodynamics states that when systems, A and B, are in thermal equilibrium with each other, and B is in
thermal equilibrium with a third system, C, then A is also in thermal equilibrium with C.
Two systems are in thermal equilibrium if they could transfer heat between each other, but don’t.
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If enough time is allowed for heat transfer to occur between a thermometer and a system, the temperature the thermometer
registers does represent the system with which it achieves thermal equilibrium.
Key Terms
heat: energy transferred from one body to another by thermal interactions
thermodynamics: a branch of natural science concerned with heat and its relation to energy and work
internal energy: The sum of all energy present in the system, including kinetic and potential energy; equivalently, the energy
needed to create a system, excluding the energy necessary to displace its surroundings.
thermal equilibrium: Two systems are in thermal equilibrium if they could transfer heat between each other, but don’t.
thermodynamic temperature: Temperature defined in terms of the laws of thermodynamics rather than the properties of a real
material: expressed in kelvins
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14.2: The First Law of Thermodynamics
learning objectives
Explain how the net heat transferred and net work done in a system relate to the first law of thermodynamics
The first law of thermodynamics is a version of the law of conservation of energy specialized for thermodynamic systems. It is
usually formulated by stating that the change in the internal energy of a closed system is equal to the amount of heat supplied to the
system, minus the amount of work done by the system on its surroundings. The law of conservation of energy can be stated like
this: The energy of an isolated system is constant.
First Law of Thermodynamics: In this video I continue with my series of tutorial videos on Thermal Physics and
Thermodynamics. It’s pitched at undergraduate level and while it is mainly aimed at physics majors, it should be useful to anybody
taking a first course in thermodynamics such as engineers etc..
If we are interested in how heat transfer is converted into work, then the conservation of energy principle is important. The first law
of thermodynamics applies the conservation of energy principle to systems where heat transfer and doing work are the methods of
transferring energy into and out of the system. In equation form, the first law of thermodynamics is
14.2.1 https://phys.libretexts.org/@go/page/14523
Internal Energy: The first law of thermodynamics is the conservation-of-energy principle stated for a system where heat and work
are the methods of transferring energy for a system in thermal equilibrium. Q represents the net heat transfer—it is the sum of all
heat transfers into and out of the system. Q is positive for net heat transfer into the system. W is the total work done on and by the
system. W is positive when more work is done by the system than on it. The change in the internal energy of the system, ΔU, is
related to heat and work by the first law of thermodynamics, ΔU = Q − W .
ΔU = Q − W. (14.2.1)
Here ΔU is the change in internal energy U of the system, Q is the net heat transferred into the system, and W is the net work done
by the system. We use the following sign conventions: if Q is positive, then there is a net heat transfer into the system; if W is
positive, then there is net work done by the system. So positive Q adds energy to the system and positive W takes energy from the
system. Thus ΔU = Q − W . Note also that if more heat transfer into the system occurs than work done, the difference is stored as
internal energy. Heat engines are a good example of this—heat transfer into them takes place so that they can do work.
learning objectives
Contrast isobaric and isochoric processes
According to the first law of thermodynamics, heat transferred to a system can be either converted to internal energy or used to do
work to the environment. A process in which a gas does work on its environment at constant pressure is called an isobaric process,
while one in which volume is kept constant is called an isochoric process.
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Fig 1: An isobaric expansion of a gas requires heat transfer during the expansion to keep the pressure constant. Since pressure is
constant, the work done is PΔV.
We may say that the system is dynamically insulated, by a rigid boundary, from the environment.
Isothermal Processes
An isothermal process is a change of a thermodynamic system, in which the temperature remains constant.
learning objectives
Identify the typical systems in which an isothermal process occurs
An isothermal process is a change of a system, in which the temperature remains constant: ΔT = 0 . This typically occurs when a
system is in contact with an outside thermal reservoir ( heat bath), and the change occurs slowly enough to allow the system to
continually adjust to the temperature of the reservoir through heat exchange. In contrast, an adiabatic process is where a system
exchanges no heat with its surroundings (Q = 0 ). (See our atom on “Adiabatic Process. “) In other words, in an isothermal process,
the value ΔT = 0 but Q ≠ 0 , while in an adiabatic process, ΔT ≠ 0 but Q = 0 .
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nRT constant
P = = (14.2.3)
V V
holds.
Calculation of Work
In thermodynamics, the work involved when a gas changes from state A to state B is simply
VB
For an isothermal, reversible process, this integral equals the area under the relevant pressure-volume isotherm, and is indicated in
blue in for an ideal gas. Again, P = nRT
V
applies and with T being constant (as this is an isothermal process), we have
Work Done by Gas During Expansion: The blue area represents “work” done by the gas during expansion for this isothermal
change.
VB
1 VB
WA→B = nRT ∫ dV = nRT ln . (14.2.5)
VA
V VA
It is also worth noting that, for many systems, if the temperature is held constant, the internal energy of the system also is constant,
and so ΔU=0ΔU=0. From the first law of thermodynamics, it follows that Q=−WQ=−W for this same isothermal process.
Adiabatic Processes
An adiabatic process is any process occurring without gain or loss of heat within a system.
learning objectives
Assess the environments in which isothermal processes typically occur
An isothermal process is a change of a system, in which the temperature remains constant: ΔT = 0 . This typically occurs when a
system is in contact with an outside thermal reservoir (heat bath), and the change occurs slowly enough to allow the system to
continually adjust to the temperature of the reservoir through heat exchange. In contrast, an adiabatic process is where a system
exchanges no heat with its surroundings (Q = 0 ). (See our atom on “Adiabatic Process. “) In other words, in an isothermal process,
the value ΔT = 0 but Q ≠ 0 , while in an adiabatic process, ΔT ≠ 0 but Q = 0 .
nRT constant
P = = (14.2.6)
V V
holds. The family of curves generated by this equation is shown in. Each curve is called an isotherm.
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Isotherms of an Ideal Gas: Several isotherms of an ideal gas on a PV diagram.
Calculation of Work
In thermodynamics, the work involved when a gas changes from state A to state B is simply
VB
For an isothermal, reversible process, this integral equals the area under the relevant pressure-volume isotherm, and is indicated in
blue in for an ideal gas. Again, P = nRT / V applies and with T being constant (as this is an isothermal process), we have
Work Done by Gas During Expansion: The blue area represents “work” done by the gas during expansion for this isothermal
change.
VB
1 VB
WA→B = nRT ∫ dV = nRT ln . (14.2.8)
VA V VA
It is also worth noting that, for many systems, if the temperature is held constant, the internal energy of the system also is constant,
and so ΔU = 0 . From the first law of thermodynamics, it follows that Q = −W for this same isothermal process.
Human Metabolism
The 1st law of thermodynamics explains human metabolism: the conversion of food into energy that is used by the body to perform
activities.
14.2.5 https://phys.libretexts.org/@go/page/14523
learning objectives
Contrast catabolism and anabolism in regards to energy
Metabolism in humans is the conversion of food into energy, which is then used by the body to perform activities. It is an example
of the first law of thermodynamics in action. Considering the body as the system of interest, we can use the first law to examine
heat transfer, doing work, and internal energy in activities ranging from sleep to heavy exercise. For example, one major factor in
such activities is body temperature—normally kept constant by heat transfer to the surroundings, meaning that Q is negative (i.e.,
our body loses heat). Another factor is that the body usually does work on the outside world, meaning that W is positive. Thus, in
such situations the body loses internal energy, since ΔU = Q − W is negative.
Eating
Now consider the effects of eating. The body metabolizes all the food we consume. Eating increases the internal energy of the body
by adding chemical potential energy. In essence, metabolism uses an oxidation process in which the chemical potential energy of
food is released. This implies that food input is in the form of work. Food energy is reported in a special unit, known as the Calorie.
This energy is measured by burning food in a calorimeter, which is how the units are determined.
Internal energy
Our body loses internal energy, and there are three places this internal energy can go—to heat transfer, to doing work, and to stored
fat (a tiny fraction also goes to cell repair and growth). As shown in Fig 1 heat transfer and doing work take internal energy out of
the body, and then food puts it back. If you eat just the right amount of food, then your average internal energy remains constant.
Whatever you lose to heat transfer and doing work is replaced by food, so that, in the long run, ΔU=0. If you overeat repeatedly,
then ΔU is always positive, and your body stores this extra internal energy as fat. The reverse is true if you eat too little. If ΔU is
negative for a few days, then the body metabolizes its own fat to maintain body temperature and do work that takes energy from the
body. This process is how dieting produces weight loss.
Metabolism: (a) The first law of thermodynamics applied to metabolism. Heat transferred out of the body (Q) and work done by
the body (W) remove internal energy, while food intake replaces it. (Food intake may be considered as work done on the body. ) (b)
Plants convert part of the radiant heat transfer in sunlight to stored chemical energy, a process called photosynthesis.
Metabolism
Life is not always this simple, as any dieter knows. The body stores fat or metabolizes it only if energy intake changes for a period
of several days. Once you have been on a major diet, the next one is less successful because your body alters the way it responds to
low energy intake. Your basal metabolic rate is the rate at which food is converted into heat transfer and work done while the body
is at complete rest. The body adjusts its basal metabolic rate to compensate (partially) for over-eating or under-eating. The body
will decrease the metabolic rate rather than eliminate its own fat to replace lost food intake. You will become more easily chilled
and feel less energetic as a result of the lower metabolic rate, and you will not lose weight as fast as before. Exercise helps with
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weight loss because it produces both heat transfer from your body and work, and raises your metabolic rate even when you are at
rest.
Irreversibility
The body provides us with an excellent indication that many thermodynamic processes are irreversible. An irreversible process can
go in one direction but not the reverse, under a given set of conditions. For example, although body fat can be converted to do work
and produce heat transfer, work done on the body and heat transfer into it cannot be converted to body fat. Otherwise, we could
skip lunch by sunning ourselves or by walking down stairs. Another example of an irreversible thermodynamic process is
photosynthesis. This process is the intake of one form of energy—light—by plants and its conversion to chemical potential energy.
Both applications of the first law of thermodynamics are illustrated in. One great advantage of such conservation laws is that they
accurately describe the beginning and ending points of complex processes (such as metabolism and photosynthesis) without regard
to the complications in between.
Key Points
The first law of thermodynamics is a version of the law of conservation of energy, specialized for thermodynamical systems.
In equation form, the first law of thermodynamics is ΔU = Q − W .
Heat engines are a good example of the application of the 1st law; heat transfer into them takes place so that they can do work.
An isobaric process occurs at constant pressure. Since the pressure is constant, the force exerted is constant and the work done
is given as PΔV.
An isobaric expansion of a gas requires heat transfer to keep the pressure constant.
An isochoric process is one in which the volume is held constant, meaning that the work done by the system will be zero. The
only change will be that a gas gains internal energy.
For an ideal gas, the product of pressure and volume (PV) is a constant if the gas is kept at isothermal conditions.
For an ideal gas, the work involved when a gas changes from state A to state B through an isothermal process is given as
VB
W A→B = nRT ln
VA
.
For many systems, if the temperature is held constant, the internal energy of the system also is constant. It follows that
Q = −W in this case.
Adiabatic processes can occur if the container of the system has thermally-insulated walls or the process happens in an
extremely short time.
For an adiabatically expanding ideal monatomic gas which does work on its environment (W is positive), internal energy of the
gas should decrease.
In a sense, isothermal process can be considered as the opposite extreme of adiabatic process. In isothermal processes, heat
exchange is slow enough so that the system’s temperature remains constant.
Human metabolism is a complicated process. The 1st law of thermodynamics describes the beginning and ending points of
these processes.
Our body loses internal energy. There are three places this internal energy can go—to heat transfer, to doing work, and to stored
fat.
Our body provides a good example of irreversible processes. Although body fat can be converted to do work and produce heat
transfer, work done on the body and heat transfer into it cannot be converted to body fat.
Key Terms
internal energy: The sum of all energy present in the system, including kinetic and potential energy; equivalently, the energy
needed to create a system, excluding the energy necessary to displace its surroundings.
heat: energy transferred from one body to another by thermal interactions
law of conservation of energy: The law stating that the total amount of energy in any isolated system remains constant, and
cannot be created or destroyed, although it may change forms.
internal energy: The sum of all energy present in the system, including kinetic and potential energy; equivalently, the energy
needed to create a system, excluding the energy necessary to displace its surroundings.
reversible: Capable of returning to the original state without consumption of free energy and increase of entropy.
ideal gas: A hypothetical gas whose molecules exhibit no interaction and undergo elastic collision with each other and with the
walls of the container.
Boyle’s law: The observation that the pressure of an ideal gas is inversely proportional to its volume at constant temperature.
14.2.7 https://phys.libretexts.org/@go/page/14523
ideal gas: A hypothetical gas whose molecules exhibit no interaction and undergo elastic collision with each other and with the
walls of the container.
metabolism: The complete set of chemical reactions that occur in living cells.
oxidation: A reaction in which the atoms of an element lose electrons and the valence of the element increases.
calorie: The energy needed to increase the temperature of 1 kilogram of water by 1 kelvin. It is equivalent to 1,000 (small)
calories.
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14.3: The Second Law of Thermodynamics
learning objectives
Contrast the concept of irreversibility between the First and Second Laws of Thermodynamics
Irreversibility
The second law of thermodynamics deals with the direction taken by spontaneous processes. Many processes occur spontaneously
in one direction only—that is, they areirreversible, under a given set of conditions. Although irreversibility is seen in day-to-day
life—a broken glass does not resume its original state, for instance—complete irreversibility is a statistical statement that cannot be
seen during the lifetime of the universe. More precisely, an irreversible process is one that depends on path. If the process can go in
only one direction, then the reverse path differs fundamentally and the process cannot be reversible.
For example, heat involves the transfer of energy from higher to lower temperature. A cold object in contact with a hot one never
gets colder, transferring heat to the hot object and making it hotter. Furthermore, mechanical energy, such as kinetic energy, can be
completely converted to thermal energy by friction, but the reverse is impossible. A hot stationary object never spontaneously cools
off and starts moving. Yet another example is the expansion of a puff of gas introduced into one corner of a vacuum chamber. The
gas expands to fill the chamber, but it never regroups in the corner. The random motion of the gas molecules could take them all
back to the corner, but this is never observed to happen.
One-Way Processed in Nature: Examples of one-way processes in nature. (a) Heat transfer occurs spontaneously from hot to cold
and not from cold to hot. (b) The brakes of this car convert its kinetic energy to heat transfer to the environment. The reverse
process is impossible. (c) The burst of gas let into this vacuum chamber quickly expands to uniformly fill every part of the
chamber. The random motions of the gas molecules will never return them to the corner.
14.3.1 https://phys.libretexts.org/@go/page/14524
AP Physics 2 - Thermodynamics and Heat Engines
Thermodynamics and Heat Engines: A brief introduction to heat engines and thermodynamic concepts such as the Carnot Engine
for students.
The Second Law of Thermodynamics(first expression): Heat transfer occurs spontaneously from higher- to lower-temperature
bodies but never spontaneously in the reverse direction.
The law states that it is impossible for any process to have as its sole result heat transfer from a cooler to a hotter object. We will
express the law in other terms later on, most importantly in terms of entropy.
Heat Engines
In thermodynamics, a heat engine is a system that performs the conversion of heat or thermal energy to mechanical work.
learning objectives
Justify why efficiency is one of the most important parameters for any heat engine
In thermodynamics, a heat engine is a system that performs the conversion of heat or thermal energy to mechanical work. Gasoline
and diesel engines, jet engines, and steam turbines are all heat engines that do work by using part of the heat transfer from some
source. Heat transfer from the hot object (or hot reservoir) is denoted as Qh, while heat transfer into the cold object (or cold
reservoir) is Qc, and the work done by the engine is W. The temperatures of the hot and cold reservoirs are Th and Tc, respectively.
14.3.2 https://phys.libretexts.org/@go/page/14524
Heat Transfer: (a) Heat transfer occurs spontaneously from a hot object to a cold one, consistent with the second law of
thermodynamics. (b) A heat engine, represented here by a circle, uses part of the heat transfer to do work. The hot and cold objects
are called the hot and cold reservoirs. Qh is the heat transfer out of the hot reservoir, W is the work output, and Qc is the heat
transfer into the cold reservoir.
14.3.3 https://phys.libretexts.org/@go/page/14524
Thermodynamics and Heat Engines: A brief introduction to heat engines and thermodynamic concepts such as the Carnot Engine
for students.
Because the hot reservoir is heated externally, which is energy intensive, it is important that the work is done as efficiently as
possible. In fact, we would like W to equal Qh, and for there to be no heat transfer to the environment (Qc=0). Unfortunately, this is
impossible. The second law of thermodynamics (second expression) also states, with regard to using heat transfer to do work: It is
impossible in any system for heat transfer from a reservoir to completely convert to work in a cyclical process in which the system
returns to its initial state.
A cyclical process brings a system, such as the gas in a cylinder, back to its original state at the end of every cycle. Most heat
engines, such as reciprocating piston engines and rotating turbines, use cyclical processes. The second law, in its second form,
clearly states that such engines cannot have perfect conversion of heat transfer into work done.
Efficiency
A cyclical process brings the system back to its original condition at the end of every cycle. By definition, such a system’s internal
energy U is the same at the beginning and end of every cycle—that is, ΔU = 0 . The first law of thermodynamics states that
ΔU = Q − W , where Q is the net heat transfer during the cycle (W = Q − Q ) and W is the net work done by the system. Since
h c
ΔU = 0 for a complete cycle, we have W = Q . Thus the net work done by the system equals the net heat transfer into the system,
or
W = Qh − Qc (cyclical process),
just as shown schematically in (b).
Efficiency is one of the most important parameters for any heat engine. The problem is that in all processes, there is significant heat
transfer Qc lost to the environment. In the conversion of energy to work, we are always faced with the problem of getting less out
than we put in. We define the efficiency of a heat engine (Eff) to be its net work output W divided by heat transfer to the engine Qh:
Ef f = W Qh . (14.3.1)
making it clear that an efficiency of 1, or 100%, is possible only if there is no heat transfer to the environment (Qc=0).
Carnot Cycles
The Carnot cycle is the most efficient cyclical process possible and uses only reversible processes through its cycle.
learning objectives
Analyze why the Carnot engine is considered the perfect engine
We know from the second law of thermodynamics that a heat engine cannot be 100 percent efficient, since there must always be
some heat transfer Qc to the environment. (See our atom on “Heat Engines. “) How efficient can a heat engine be then? This
question was answered at a theoretical level in 1824 by a young French engineer, Sadi Carnot (1796-1832), in his study of the then-
emerging heat engine technology crucial to the Industrial Revolution. He devised a theoretical cycle, now called the Carnot cycle,
which is the most efficient cyclical process possible. The second law of thermodynamics can be restated in terms of the Carnot
cycle, and so what Carnot actually discovered was this fundamental law. Any heat engine employing the Carnot cycle is called a
Carnot engine.
What is crucial to the Carnot cycle is that only reversible processes are used. Irreversible processes involve dissipative factors, such
as friction and turbulence. This increases heat transfer Qc to the environment and reduces the efficiency of the engine. Obviously,
then, reversible processes are superior.
The second law of thermodynamics (a third form): A Carnot engine operating between two given temperatures has the greatest
possible efficiency of any heat engine operating between these two temperatures. Furthermore, all engines employing only
reversible processes have this same maximum efficiency when operating between the same given temperatures.
14.3.4 https://phys.libretexts.org/@go/page/14524
Efficiency
The Carnot cycle comprises two isothermal and two adiabatic processes. Recall that both isothermal and adiabatic processes are, in
principle, reversible.
PV Diagram for a Carnot Cycle: PV diagram for a Carnot cycle, employing only reversible isothermal and adiabatic processes.
Heat transfer Qh occurs into the working substance during the isothermal path AB, which takes place at constant temperature Th.
Heat transfer Qc occurs out of the working substance during the isothermal path CD, which takes place at constant temperature Tc.
The net work output W equals the area inside the path ABCDA. Also shown is a schematic of a Carnot engine operating between
hot and cold reservoirs at temperatures Th and Tc.
Carnot also determined the efficiency of a perfect heat engine—that is, a Carnot engine. It is always true that the efficiency of a
Qh −Qc Qc
cyclical heat engine is given by: Ef f = Qh
=1−
Qh
.
What Carnot found was that for a perfect heat engine, the ratio \mathrm{\frac{Q_c}{Q_h}}} equals the ratio of the absolute
Qc Tc
temperatures of the heat reservoirs. That is, Qh
=
Th
for a Carnot engine, so that the maximum or Carnot efficiency EffC is given
Tc
by Ef f = 1 −
c
Th
, where Th and Tc are in kelvins. (Derivation of the formula is slightly beyond the scope of this atom. ) No real
heat engine can do as well as the Carnot efficiency—an actual efficiency of about 0.7 of this maximum is usually the best that can
be accomplished.
learning objectives
Explain how the the components of a heat pump cause heat to transfer from a cold reservoir to a hot reservoir
Heat pumps, air conditioners, and refrigerators utilize heat transfer from cold to hot. Heat transfer (Qc) occurs from a cold reservoir
and into a hot one. This requires work input W, which is also converted to heat transfer. Thus the heat transfer to the hot reservoir is
Qh=Qc+W. A heat pump’s mission is for heat transfer Qh to occur into a warm environment, such as a home in the winter. The
mission of air conditioners and refrigerators is for heat transfer Qc to occur from a cool environment, such as chilling a room or
keeping food at lower temperatures than the environment. Actually, a heat pump can be used both to heat and cool a space. It is
essentially an air conditioner and a heating unit all in one. In this section we will concentrate on its heating mode.
Heat Pumps
A working fluid such as a non-CFC refrigerant is used in a basic heat pump. The basic components of a heat pump in are a
condenser, an expansion valve, an evaporator and a compressor. In the outdoor coils (the evaporator), heat transfer Qc occurs to the
working fluid from the cold outdoor air, turning it into a gas. The electrically driven compressor (work input W) raises the
temperature and pressure of the gas and forces it into the condenser coils that are inside the heated space. Because the temperature
of the gas is higher than the temperature inside the room, heat transfer to the room occurs and the gas condenses to a liquid. The
liquid then flows back through a pressure-reducing valve to the outdoor evaporator coils, being cooled through expansion. (In a
cooling cycle, the evaporator and condenser coils exchange roles and the flow direction of the fluid is reversed. )
14.3.5 https://phys.libretexts.org/@go/page/14524
Simple Heat Pump: A simple heat pump has four basic components: (1) condenser, (2) expansion valve, (3) evaporator, and (4)
compressor.
Coefficient of Performance
The quality of a heat pump is judged by how much heat transfer Qhoccurs into the warm space compared with how much work
input W is required. We define a heat pump’s coefficient of performance (COPhp) to be
Qh
COPhp = . (14.3.2)
W
Since the efficiency of a heat engine is Eff=W/Qh, we see that COPhp=1/Eff. Since the efficiency of any heat engine is less than 1, it
means that COPhp is always greater than 1—that is, a heat pump always has more heat transfer Qh than work put into it. Another
interesting point is that heat pumps work best when temperature differences are small. The efficiency of a perfect engine (or Carnot
engine) is
Tc
Ef fC = 1 ; (14.3.3)
Th
thus, the smaller the temperature difference, the smaller the efficiency and the greater the COPhp.
Qh
Since Q h = Qc + W and COP hp =
W
, we derive that
Also, from Qh>Qc, we see that an air conditioner will have a lower coefficient of performance than a heat pump.
Key Points
Many thermodynamic phenomena, allowed to occur by the first law of thermodynamics, never occur in nature.
Many processes occur spontaneously in one direction only, and the second law of thermodynamics deals with the direction
taken by spontaneous processes.
According to the second law of thermodynamics, it is impossible for any process to have heat transfer from a cooler to a hotter
object as its sole result.
14.3.6 https://phys.libretexts.org/@go/page/14524
A cyclical process brings a system, such as the gas in a cylinder, back to its original state at the end of every cycle. Most heat
engines, such as reciprocating piston engines and rotating turbines, use cyclical processes.
The second law of thermodynamics can be expressed as the following: It is impossible in any system for heat transfer from a
reservoir to completely convert to work in a cyclical process in which the system returns to its initial state.
The efficiency of a heat engine (Eff) is defined to be the engine’s net work output W divided by heat transfer to the engine:
Eff=WQh=1−QcQhEff=WQh=1−QcQh, where Qc and Qh denotes heat transfer to hot (engine) and cold (environment)
reservoir.
The second law of thermodynamics indicates that a Carnot engine operating between two given temperatures has the greatest
possible efficiency of any heat engine operating between these two temperatures.
Irreversible processes involve dissipative factors, which reduces the efficiency of the engine. Obviously, reversible processes
are superior from the efficiency perspective.
Carnot efficiency, the maximum achievable heat engine efficiency, is given as Effc=1−TcThEffc=1−TcTh.
A heat pump ‘s mission is for heat transfer Qh to occur into a warm environment, such as a home in the winter.
The mission of air conditioners and refrigerators is for heat transfer Qc to occur from a cool environment, such as chilling a
room or keeping food at lower temperatures than the environment.
A heat pump can be used both to heat and cool a space. It is essentially an air conditioner and a heating unit all in one. This is
made possible by reversing the flow of its refrigerant, changing the direction net heat transfer.
Key Terms
entropy: A measure of how evenly energy (or some analogous property) is distributed in a system.
the first law of thermodynamics: A version of the law of conservation of energy, specialized for thermodynamical systems.
Usually expressed as ΔU=Q−W.
thermal energy: The internal energy of a system in thermodynamic equilibrium due to its temperature.
internal energy: The sum of all energy present in the system, including kinetic and potential energy; equivalently, the energy
needed to create a system, excluding the energy necessary to displace its surroundings.
the second law of thermodynamics: A law stating that states that the entropy of an isolated system never decreases, because
isolated systems spontaneously evolve toward thermodynamic equilibrium—the state of maximum entropy. Equivalently,
perpetual motion machines of the second kind are impossible.
heat engine: Any device which converts heat energy into mechanical work.
CFC: An organic compound that was commonly used as a refrigerant. Not commonly used anymore because of its ozone
depletion effect.
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14.4: Entropy
learning objectives
Calculate the total change in entropy for a system in a reversible process
In this and following Atoms, we will study entropy. By examining it, we shall see that the directions associated with the second law
— heat transfer from hot to cold, for example—are related to the tendency in nature for systems to become disordered and for less
energy to be available for use as work. The entropy of a system can in fact be shown to be a measure of its disorder and of the
unavailability of energy to do work.
Definition of Entropy
We can see how entropy is defined by recalling our discussion of the Carnot engine. We noted that for a Carnot cycle, and hence for
Qc Qh
any reversible processes, Qc/Qh=Tc/Th. Rearranging terms yields Tc
=
Th
for any reversible process. Qc and Qhare absolute
values of the heat transfer at temperatures Tc and Th, respectively. This ratio of Q/T is defined to be the change in entropy ΔS for a
reversible process,
Q
ΔS = ( )rev , (14.4.1)
T
where Q is the heat transfer, which is positive for heat transfer into and negative for heat transfer out of, and T is the absolute
temperature at which the reversible process takes place. The SI unit for entropy is joules per kelvin (J/K). If temperature changes
during the process, then it is usually a good approximation (for small changes in temperature) to take T to be the average
temperature, avoiding the need to use integral calculus to find ΔS.
The definition of ΔS is strictly valid only for reversible processes, such as used in a Carnot engine. However, we can find ΔS
precisely even for real, irreversible processes. The reason is that the entropy S of a system, like internal energy U, depends only on
the state of the system and not how it reached that condition. Entropy is a property of state. Thus the change in entropy ΔS of a
system between state one and state two is the same no matter how the change occurs. We just need to find or imagine a reversible
process that takes us from state one to state two and calculate ΔS for that process. That will be the change in entropy for any
process going from state one to state two.
Change in Entropy: When a system goes from state one to state two, its entropy changes by the same amount ΔS, whether a
hypothetical reversible path is followed or a real irreversible path is taken.
Example
Now let us take a look at the change in entropy of a Carnot engine and its heat reservoirs for one full cycle. The hot reservoir has a
loss of entropy ΔSh=−Qh/Th, because heat transfer occurs out of it (remember that when heat transfers out, then Q has a negative
sign). The cold reservoir has a gain of entropy ΔSc=Qc/Tc, because heat transfer occurs into it. (We assume the reservoirs are
sufficiently large that their temperatures are constant. ) So the total change in entropy is
14.4.1 https://phys.libretexts.org/@go/page/14525
PV Diagram for a Carnot Cycle: PV diagram for a Carnot cycle, employing only reversible isothermal and adiabatic processes.
Heat transfer Qh occurs into the working substance during the isothermal path AB, which takes place at constant temperature Th.
Heat transfer Qc occurs out of the working substance during the isothermal path CD, which takes place at constant temperature Tc.
The net work output W equals the area inside the path ABCDA. Also shown is a schematic of a Carnot engine operating between
hot and cold reservoirs at temperatures Th and Tc.
This result, which has general validity, means that the total change in entropy for a system in any reversible process is zero.
learning objectives
Calculate the number of microstates for simple configurations
The various ways of formulating the second law of thermodynamics tell what happens rather than why it happens. Why should heat
transfer occur only from hot to cold? Why should the universe become increasingly disorderly? The answer is that it is a matter of
overwhelming probability. Disorder is simply vastly more likely than order. To illustrate this fact, we will examine some random
processes, starting with coin tosses.
Coin Tosses
What are the possible outcomes of tossing 5 coins? Each coin can land either heads or tails. On the large scale, we are concerned
only with the total heads and tails and not with the order in which heads and tails appear. The following table shows all possibilities
along with numbers of possible configurations (or microstate; a detailed description of every element of a system). For example, 4
heads and 1 tail instance may occur on 5 different configurations, with any one of the 5 coins showing tail and all the rest heads.
(HHHHT, HHHTH, HHTHH, HTHHH, THHHH)
5 heads, 0 tails: 1 microstate
4 heads, 1 tail: 5 microstates
3 heads, 2 tails: 10 microstates
2 heads, 3 tails: 10 microstates
1 head, 4 tails: 5 microstates
0 head, 5 tails: 1 microstate
Note that all of these conclusions are based on the crucial assumption that each microstate is equally probable. Otherwise, the
analysis will be erroneous.
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The two most orderly possibilities are 5 heads or 5 tails. (They are more structured than the others. ) They are also the least likely,
only 2 out of 32 possibilities. The most disorderly possibilities are 3 heads and 2 tails and its reverse. (They are the least structured.
) The most disorderly possibilities are also the most likely, with 20 out of 32 possibilities for the 3 heads and 2 tails and its reverse.
If we start with an orderly array like 5 heads and toss the coins, it is very likely that we will get a less orderly array as a result, since
30 out of the 32 possibilities are less orderly. So even if you start with an orderly state, there is a strong tendency to go from order
to disorder, from low entropy to high entropy. The reverse can happen, but it is unlikely.
This result becomes dramatic for larger systems. Consider what happens if you have 100 coins instead of just 5. The most orderly
arrangements (most structured) are 100 heads or 100 tails. The least orderly (least structured) is that of 50 heads and 50 tails. There
is only 1 way (1 microstate) to get the most orderly arrangement of 100 heads. The total number of different ways 100 coins can be
tossed—is an impressively large 1.27×1030. Now, if we start with an orderly macrostate like 100 heads and toss the coins, there is a
virtual certainty that we will get a less orderly macrostate. If you tossed the coins once each second, you could expect to get either
100 heads or 100 tails once in 2×1022 years! In contrast, there is an 8% chance of getting 50 heads, a 73% chance of getting from
45 to 55 heads, and a 96% chance of getting from 40 to 60 heads. Disorder is highly likely.
Real Gas
The fantastic growth in the odds favoring disorder that we see in going from 5 to 100 coins continues as the number of entities in
the system increases. In a volume of 1 m3, roughly 1023 molecules (or the order of magnitude of Avogadro’s number) are present in
a gas. The most likely conditions (or macrostate) for the gas are those we see all the time—a random distribution of atoms in space
with a Maxwell-Boltzmann distribution of speeds in random directions, as predicted by kinetic theory as shown in (a). This is the
most disorderly and least structured condition we can imagine.
Kinetic Theory: (a) The ordinary state of gas in a container is a disorderly, random distribution of atoms or molecules with a
Maxwell-Boltzmann distribution of speeds. It is so unlikely that these atoms or molecules would ever end up in one corner of the
container that it might as well be impossible. (b) With energy transfer, the gas can be forced into one corner and its entropy greatly
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reduced. But left alone, it will spontaneously increase its entropy and return to the normal conditions, because they are immensely
more likely.
In contrast, one type of very orderly and structured macrostate has all of the atoms in one corner of a container with identical
velocities. There are very few ways to accomplish this (very few microstates corresponding to it), and so it is exceedingly unlikely
ever to occur. (See (b). ) Indeed, it is so unlikely that we have a law saying that it is impossible, which has never been observed to
be violated—the second law of thermodynamics.
Order to Disorder
Entropy is a measure of disorder, so increased entropy means more disorder in the system.
learning objectives
Discuss entropy and disorder within a system
Entropy is a measure of disorder. This notion was initially postulated by Ludwig Boltzmann in the 1800s. For example, melting a
block of ice means taking a highly structured and orderly system of water molecules and converting it into a disorderly liquid in
which molecules have no fixed positions. There is a large increase in entropy in the process.
Entropy of Ice: When ice melts, it becomes more disordered and less structured. The systematic arrangement of molecules in a
crystal structure is replaced by a more random and less orderly movement of molecules without fixed locations or orientations. Its
entropy increases because heat transfer occurs into it. Entropy is a measure of disorder.
Example 14.4.1:
As an example, suppose we mix equal masses of water originally at two different temperatures, say 20.0º C and 40.0º C. The
result is water at an intermediate temperature of 30.0º C. Three outcomes have resulted:
Entropy has increased.
Some energy has become unavailable to do work.
The system has become less orderly.
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Heat Death
The entropy of the universe is constantly increasing and is destined for thermodynamic equilibrium, called the heat death of the
universe.
learning objectives
Describe processes that lead to the heat death of the universe
In the early, energetic universe, all matter and energy were easily interchangeable and identical in nature. Gravity played a vital role
in the young universe. Although it may have seemed disorderly, there was enormous potential energy available to do work—all the
future energy in the universe.
As the universe matured, temperature differences arose, which created more opportunity for work. Stars are hotter than planets, for
example, which are warmer than icy asteroids, which are warmer still than the vacuum of the space between them. Most of these
are cooling down from their usually violent births, at which time they were provided with energy of their own—nuclear energy in
the case of stars, volcanic energy on Earth and other planets, and so on. Without additional energy input, however, their days are
numbered.
Infant Universe: The image of an infant universe reveals temperature fluctuations (shown as color differences) that correspond to
the seeds that grew to become the galaxies.
As entropy increases, less and less energy in the universe is available to do work. On Earth, we still have great stores of energy
such as fossil and nuclear fuels; large-scale temperature differences, which can provide wind energy; geothermal energies due to
differences in temperature in Earth’s layers; and tidal energies owing to our abundance of liquid water. As these are used, a certain
fraction of the energy they contain can never be converted into doing work. Eventually, all fuels will be exhausted, all temperatures
will equalize, and it will be impossible for heat engines to function, or for work to be done.
Since the universe is a closed system, the entropy of the universe is constantly increasing, and so the availability of energy to do
work is constantly decreasing. Eventually, when all stars have died, all forms of potential energy have been utilized, and all
temperatures have equalized (depending on the mass of the universe, either at a very high temperature following a universal
contraction, or a very low one, just before all activity ceases) there will be no possibility of doing work.
Either way, the universe is destined for thermodynamic equilibrium—maximum entropy. This is often called the heat death of the
universe, and will mean the end of all activity. However, whether the universe contracts and heats up, or continues to expand and
cools down, the end is not near. Calculations of black holes suggest that entropy can easily continue for at least 10100 years.
learning objectives
Formulate conditions that allow decrease of the entropy in one part of the universe
Some people misunderstand the second law of thermodynamics, stated in terms of entropy, to say that the process of the evolution
of life violates this law. Over time, complex organisms evolved from much simpler ancestors, representing a large decrease in
entropy of the Earth’s biosphere. It is a fact that living organisms have evolved to be highly structured, and much lower in entropy
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than the substances from which they grow. But it is always possible for the entropy of one part of the universe to decrease,
provided the total change in entropy of the universe increases. In equation form, we can write this as
ΔStot = ΔSsys + ΔSenv > 0. (14.4.4)
Thus ΔSsys can be negative as long as ΔSenv is positive and greater in magnitude.
How is it possible for a system to decrease its entropy? Energy transfer is necessary. If I gather iron ore from the ground and
convert it into steel and build a bridge, my work (and used energy) has decreased the entropy of that system. Energy coming from
the Sun can decrease the entropy of local systems on Earth—that is, ΔSsys is negative. But the overall entropy of the rest of the
universe increases by a greater amount—that is, ΔSenv is positive and greater in magnitude. Thus, ΔStot>0, and the second law of
thermodynamics is not violated.
Every time a plant stores some solar energy in the form of chemical potential energy, or an updraft of warm air lifts a soaring bird,
the Earth can be viewed as a heat engine operating between a hot reservoir supplied by the Sun and a cold reservoir supplied by
dark outer space—a heat engine of high complexity, causing local decreases in entropy as it uses part of the heat transfer from the
Sun into deep space. However, there is a large total increase in entropy resulting from this massive heat transfer. A small part of
this heat transfer is stored in structured systems on Earth, producing much smaller local decreases in entropy.
Earth’s Entropy: Earth’s entropy may decrease in the process of intercepting a small part of the heat transfer from the Sun into
deep space. Entropy for the entire process increases greatly while Earth becomes more structured with living systems and stored
energy in various forms.
learning objectives
Describe effect of the heat dumped into the environment on the Earth’s atmospheric temperature
The Second Law of Thermodynamics may help provide explanation for why there have been increases in Earth’s temperatures over
the last 250 years (often called “Global Warming”), and many professionals are concerned that the entropy increase of the universe
is a real threat to the environment.
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Global Land-Ocean Temperature: Global mean land-ocean temperature change from 1880 – 2012, relative to the 1951 – 1980
mean. The black line is the annual mean and the red line is the five-year running mean. The green bars show uncertainty estimates.
As an engine operates, heat flows from a heat tank of greater temperature to a heat sink of lesser temperature. In between these
states, the heat flow is turned into useful energy with the help of heat engines. As these engines operate, however, a great deal of
heat is lost to the environment due to inefficiencies. In a Carnot engine, which is the most efficient theoretical engine based on
Carnot cycle, the maximum efficiency is equal to one minus the temperature of the heat sink (Tc) divided by the temperature of the
heat source (Th).
Tc
(Ef fc = 1 − ). (14.4.5)
Th
This ratio shows that for a greater efficiency to be achieved there needs to be the greatest difference in temperature available. This
brings up two important points: optimized heat sinks are at absolute zero, and the longer engines dump heat into an isolated system
the less efficient engines will become.
Unfortunately for engine efficiency, day-to-day life never operates in absolute zero. In an average car engine, only 14% to 26% of
the fuel which is put in is actually used to make the car move forward. This means that 74% to 86% is lost heat or used to power
accessories. According to the U.S. Department of Energy, 70% to 72% of heat produced by burning fuel is heat lost by the engine.
The excess heat lost by the engine is then released into the heat sink, which in the case of many modern engines would be the
Earth’s atmosphere. As more heat is dumped into the environment, Earth’s atmospheric (or heat sink) temperature will increase.
With the entropy of the environment constantly increasing, searching for new, more efficient technologies and new non-heat
engines has become a priority.
Thermal Pollution
Thermal pollution is the degradation of water quality by any process that changes ambient water temperature.
learning objectives
Identify factors that lead to thermal pollution and its ecological effects
Thermal pollution is the degradation of water quality by any process that changes ambient water temperature. A common cause of
thermal pollution is the use of water as a coolant, for example, by power plants and industrial manufacturers. When water used as a
coolant is returned to the natural environment at a higher temperature, the change in temperature decreases oxygen supply, and
affects ecosystem composition.
As we learned in our Atom on “Heat Engines”, all heat engines require heat transfer, achieved by providing (and maintaining)
temperature difference between engine’s heat source and heat sink. Water, with its high heat capacity, works extremely well as a
coolant. But this means that cooling water should be constantly replenished to maintain its cooling capacity.
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Cooling Tower: This is a cooling tower at Gustav Knepper Power Station, Dortmund, Germany. Cooling water is circulated inside
the tower.
Ecological Effects
Elevated water temperature typically decreases the level of dissolved oxygen of water. This can harm aquatic animals such as fish,
amphibians, and other aquatic organisms. An increased metabolic rate may result in fewer resources; the more adapted organisms
moving in may have an advantage over organisms that are not used to the warmer temperature. As a result, food chains of the old
and new environments may be compromised. Some fish species will avoid stream segments or coastal areas adjacent to a thermal
discharge. Biodiversity can decrease as a result. Many aquatic species will also fail to reproduce at elevated temperatures.
Some may assume that by cooling the heated water, we can possibly fix the issue of thermal pollution. However, as we noted in our
previous Atom on “Heat Pumps and Refrigerators”, work required for the additional cooling leads to more heat exhaust into the
environment. Therefore, it makes the situation even worse.
Key Points
This ratio of Q/T is defined to be the change in entropy ΔS for a reversible process: ΔS=(QT)revΔS=(QT)rev.
Entropy is a property of state. Therefore, the change in entropy ΔS of a system between two states is the same no matter how
the change occurs.
The total change in entropy for a system in any reversible process is zero.
Each microstate is equally probable in the example of coin toss. However, as a macrostate, there is a strong tendency for the
most disordered state to occur.
When tossing 100 coins, if the coins are tossed once each second, you could expect to get either all 100 heads or all 100 tails
once in 2×1022 years.
Molecules in a gas follow the Maxwell-Boltzmann distribution of speeds in random directions, which is the most disorderly and
least structured condition out of all the possibilities.
Mixing two systems may decrease the entropy of one system, but increase the entropy of the other system by a greater amount,
producing an overall increase in entropy.
After mixing water at two different temperatures, the energy in the system that could have been used to run a heat engine is now
unavailable to do work. Also, the process made the whole system more less structured.
Entropy, unavailability of energy, and disorder are not only related but are in fact essentially equivalent.
In the early, energetic universe, all matter and energy were easily interchangeable and identical in nature.
As entropy increases, less and less energy in the universe is available to do work.
The universe is destined for thermodynamic equilibrium —maximum entropy. This is often called the heat death of the
universe, and will mean the end of all activity.
Living organisms have evolved to be highly structured, and much lower in entropy than the substances from which they grow.
It possible for a system to decrease its entropy provided the total change in entropy of the universe increases:
ΔStot=ΔSsys+ΔSenv>0ΔStot=ΔSsys+ΔSenv>0.
The Earth can be viewed as a heat engine operating between a hot reservoir supplied by the Sun and a cold reservoir supplied
by dark outer space.
As heat engines operate, a great deal of heat is lost to the environment due to inefficiencies.
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Even in a Carnot engine, which is the most efficient theoretical engine, there is a heat loss determined by the ratio of
temperature of the engine and its environment.
As more heat is dumped into the environment, Earth’s atmospheric temperature will increase.
All heat engines require heat transfer, achieved by providing (and maintaining) temperature difference between engine’s heat
source and heat sink. Cooling water is typically used to maintain the temperature difference.
Elevated water temperature typically decreases the level of dissolved oxygen of water, affecting ecosystem composition.
Cooling heated water is not a solution for thermal pollution because extra work is required for the cooling, leading to more heat
exhaust into the environment.
Key Terms
Carnot cycle: A theoretical thermodynamic cycle. It is the most efficient cycle for converting a given amount of thermal energy
into work.
reversible: Capable of returning to the original state without consumption of free energy and increase of entropy.
disorder: Absence of some symmetry or correlation in a many-particle system.
Maxwell-Boltzmann distribution: A distribution describing particle speeds in gases, where the particles move freely without
interacting with one another, except for very brief elastic collisions in which they may exchange momentum and kinetic energy.
entropy: A measure of how evenly energy (or some analogous property) is distributed in a system.
geothermal: Pertaining to heat energy extracted from reservoirs in the Earth’s interior.
asteroid: A naturally occurring solid object, which is smaller than a planet and is not a comet, that orbits a star.
Carnot cycle: A theoretical thermodynamic cycle. It is the most efficient cycle for converting a given amount of thermal energy
into work.
absolute zero: The coldest possible temperature: zero on the Kelvin scale and approximately -273.15°C and -459.67°F. The
total absence of heat; the temperature at which motion of all molecules would cease.
heat engine: Any device which converts heat energy into mechanical work.
heat pump: A device that transfers heat from something at a lower temperature to something at a higher temperature by doing
work.
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14.5: The Third Law of Thermodynamics
learning objectives
Explain how absolute zero affects entropy
The third law of thermodynamics is sometimes stated as follows: The entropy of a perfect crystal at absolute zero is exactly equal
to zero.
At zero kelvin the system must be in a state with the minimum possible energy, thus this statement of the third law holds true if the
perfect crystal has only one minimum energy state. Entropy is related to the number of possible microstates, and with only one
microstate available at zero kelvin the entropy is exactly zero.
The third law was developed by the chemist Walther Nernst during the years 1906-1912. It is often referred to as Nernst’s theorem
or Nernst’s postulate. Nernst proposed that the entropy of a system at absolute zero would be a well-defined constant. Instead of
being 0, entropy at absolute zero could be a nonzero constant, due to the fact that a system may have degeneracy (having several
ground states at the same energy).
In simple terms, the third law states that the entropy of a perfect crystal approaches zero as the absolute temperature approaches
zero. This law provides an absolute reference point for the determination of entropy. ( diagrams the temperature entropy of
nitrogen. ) The entropy (S) determined relative to this point is the absolute entropy represented as follows:
Temperature Entropy of Nitrogen: Temperature–entropy diagram of nitrogen. The red curve at the left is the melting curve.
Absolute value of entropy can be determined shown here, thanks to the third law of thermodynamics.
S = kB log W, (14.5.1)
where kB is the Bolzmann constant and W is the number of microstates. Provided that the ground state is unique (or W = 1 ), the
entropy of a perfect crystal lattice as defined by Nernst’s theorem is zero provided that its ground state is unique, because
log(1) = 0.
Adiabatic Processes
It is impossible to reduce the temperature of any system to zero temperature in a finite number of finite operations.
learning objectives
Illustrate isentropic process, for example in terms of adiabtic demangetization
In our Atom on “Adiabatic Processes” (category: the First Law of Thermodynamics), we learned that an adiabatic process is any
process occurring without gain or loss of heat within a system. We also learned a monatomic ideal gas expands adiabatically. In this
14.5.1 https://phys.libretexts.org/@go/page/14526
Atom, we discuss an adiabatic cooling process that can be used to cool a gas, as well as whether absolute zero can be obtained in
real systems.
Absolute Zero?
Previously, we learned about the third law of thermodynamics, which states: the entropy of a perfect crystal at absolute zero is
exactly equal to zero.
According to the third law, the reason that T=0 cannot be reached is explained as follows: Suppose the temperature of a substance
can be reduced in an isentropic process by changing the parameter X from X to X . As an example, one can think of a multistage
2 1
adiabatic magnetization-demagnetization cycle setup where a magnetic field is switched on and off in a controlled way. (See below.
The parameter X in this case would be the magnetization of the gas. )
Assuming an entropy difference at absolute zero, T=0 could be reached in a finite number of steps. However, going back to the
third law, at T=0 there is no entropy difference, and therefore an infinite number of stepswould be needed for this process
(illustrated in ).
Can Absolute Zero be Reached?: Temperature-Entropy diagram. Horizontal lines represent isentropic processes, while vertical
lines represent isothermal processes. Left side: Absolute zero can be reached in a finite number of steps if S(T=0,X1)≠S(T=0, X2).
Right: An infinite number of steps is needed since S(0,X1)= S(0,X2).
Key Points
Entropy is related to the number of possible microstates, and with only one microstate available at zero kelvin, the entropy is
exactly zero.
The third law of thermodynamics provides an absolute reference point for the determination of entropy. The entropy determined
relative to this point is the absolute entropy.
Absolute entropy can be written as S = k \logW, where W is the number of available microstates.
B
Since at T = 0 there is no entropy difference, an infinite number of steps in a thermodynamic cooling process is required to
reach T=0.
In adiabatic demagnetization cooling, energy transfers from thermal entropy to magnetic entropy (or disorder of the magnetic
moments).
The fact that T=0 cannot be achieved in reality is a direct consequence of the third law of thermodynamics.
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Key Terms
microstate: The specific detailed microscopic configuration of a system.
absolute zero: The coldest possible temperature: zero on the Kelvin scale and approximately -273.15°C and -459.67°F. The
total absence of heat; the temperature at which motion of all molecules would cease.
degeneracy: Two or more different quantum states are said to be degenerate if they are all at the same energy level.
isentropic: Having a constant entropy.
demagnetization: The process of removing the magnetic field from an object.
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CHAPTER OVERVIEW
This page titled 15: Waves and Vibrations is shared under a not declared license and was authored, remixed, and/or curated by Boundless.
1
15.1: Introduction
learning objectives
Contrast mechanical and electromagnetic waves
Overview
A wave is an oscillation that travels through space, accompanied by a transfer of energy. Wave motion transfers energy from one
point to another, often with no permanent displacement of the particles of the medium —that is, with little or no associated mass
transport. They consist, instead, of oscillations or vibrations around almost fixed locations. There are two main types of waves.
Mechanical waves propagate through a medium, and the substance of this medium is deformed. The deformation reverses itself
owing to restoring forces resulting from its deformation.
The second main type of wave, electromagnetic waves, do not require a medium (although they may still propagate through a
medium). Instead, they consist of periodic oscillations in electrical and magnetic fields generated by charged particles, and can
therefore travel through a vacuum.
Wave Equation
The shape of a wave can take the form of any function that repeats itself over some characteristic spatial scale λλ, the wavelength
2
(see ). More generally, waveforms are scalar functions uu which satisfy the wave equation, ∂ u
2
2
=c ∇ u
2
. This equation simply
∂t
states that the acceleration of the waveform (Left: second derivative with respect to time) is proportional to the Laplacian (Right:
second spatial derivative) of the same waveform. The constant of proportionality, c , is the square of the propagation speed of the
2
wave.
Common waveforms: A sample of several common, simple waveforms. A waveform is a function that repeats in space.
Sine Wave
Plot of Sine: The sine function graphed on the Cartesian plane. In this graph, the angle x is given in radians (π = 180°).
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Consider one of the most common waveforms, the sinusoid. A general form of a sinusoidal wave is
y(x, t) = A sin(kx − ωt + ϕ) , where A is the amplitude of the wave, ω is the wave’s angular frequency, k is the wavenumber, and
ϕ is the phase of the sine wave given in radians. This waveform gives the displacement position (“y”) of a particle in a medium
from its equilibrium as a function of both position “x” and time “t”.
By taking derivatives, it is evident that the wave equation given above holds for c = , which is also called the phase speed of the
ω
wave. To find the velocity of a particle in the medium at x and t, we take the temporal derivative of the waveform to get
∂y(x, t)∂t = −Aω cos(kx − ωt + ϕ) . Likewise, to find the acceleration of the displaced particle in the medium at x and t, we
2
∂ y(x,t)
take the second derivative to get = −Aω sin(kx − ωt + ϕ) . Note the phase relationship among the trigonometric
2
∂t
2
functions in y(x, t), y (x, t), y”(x, t) . When the particle displacement is maximum or minimum, the velocity is 0. When the
′
displacement is 0, particle velocity is either maximum or minimum. Similarly, the particle acceleration is maximum (or minimum)
when the particle displacement is minimum (or maximum), respectively.
Arbitrary Wave
We looked closely into the sinusoidal wave. But how about waves that has a general form? One important aspect of the wave
equation is its linearity: the wave equation is linear in u and it is left unaltered by translations in space and time. Since a wave with
an arbitrary shape can be represented by a sum of many sinusoidal waves (this is called Fourier analysis), we can generate a great
variety of solutions of the wave equation by translating and summing sine waves that we just looked closely into.
Key Points
The wave equation requires that the second time derivative of the waveform be proportional to its second spatial derivative.
Waveforms describe the shape of physical waves, and can take the form of any function that repeats in space.
One of the most common waveforms in physics is the sinusoid. Since any arbitray waveform can be generated by adding a set
of sine waves, the physics governing a wave of an arbitrary shape can be described by using its sine wave components.
Key Terms
waveform: The shape of a physical wave, such as sound or electromagnetic radiation. The shape can be any function that
repeats in space.
Fourier analysis: The study of the way general functions may be represented or approximated by sums of simpler
trigonometric functions.
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15.2: Hooke’s Law
learning objectives
Generate the mathematical expression of Hooke’s law
In mechanics (physics), Hooke’s law is an approximation of the response of elastic (i.e., springlike) bodies. It states: the extension
of a spring is in direct proportion with the load applied to it. For instance, the spring is pulled downwards with either no load, Fp, or
twice Fp.
Diagram of Hooke’s Law: The extension of the spring is linearly proportional to the force.
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Springs and Hooke’s Law: A brief overview of springs, Hooke’s Law, and elastic potential energy for algebra-based physics
students.
Many materials obey this law of elasticity as long as the load does not exceed the material’s elastic limit. Materials for which
Hooke’s law is a useful approximation are known as linear-elastic or “Hookean” materials. Hookean materials are broadly defined
and include springs as well as muscular layers of the heart. In simple terms, Hooke’s law says that stress is directly proportional to
strain. Mathematically, Hooke’s law is stated as:
F = −kx (15.2.1)
where:
x is the displacement of the spring’s end from its equilibrium position (a distance, in SI units: meters);
F is the restoring force exerted by the spring on that end (in SI units: N or kg·m/s2); and
k is a constant called the rate or spring constant (in SI units: N/m or kg/s2). When this holds, the behavior is said to be linear. If
shown on a graph, the line should show a direct variation.
It’s possible for multiple springs to act on the same point. In such a case, Hooke’s law can still be applied. As with any other set of
forces, the forces of many springs can be combined into one resultant force.
When Hooke’s law holds, the behavior is linear; if shown on a graph, the line depicting force as a function of displacement should
show a direct variation. There is a negative sign on the right hand side of the equation because the restoring force always acts in the
opposite direction of the displacement (for example, when a spring is stretched to the left, it pulls back to the right).
Hooke’s law is named after the 17th century British physicist Robert Hooke, and was first stated in 1660 as a Latin anagram, whose
solution Hooke published in 1678 as Ut tensio, sic vis, meaning, “As the extension, so the force.”
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Hooke’s Law: The red line in this graph illustrates how force, F, varies with position according to Hooke’s law. The slope of this
line corresponds to the spring constant k. The dotted line shows what the actual (experimental) plot of force might look like. The
pictures of spring states at the bottom of the graph correspond to some points of the plot; the middle one is in the relaxed state (no
force applied).
learning objectives
Express elastic energy stored in a spring in a mathematical form
the applied force versus deformation x for a system that can be described by Hooke’s law. Work done on the system is force
multiplied by distance, which equals the area under the curve, or kx (Method A in the figure). Another way to determine the
1
2
2
work is to note that the force increases linearly from 0 to kx , so that the average force is kx , the distance moved is x, and thus
1
Applied force versus deformation: A graph of applied force versus distance for the deformation of a system that can be described
by Hooke’s law is displayed. The work done on the system equals the area under the graph or the area of the triangle, which is half
its base multiplied by its height, or W = kx . 1
2
2
W = Fapp d = (
1
2
kx)(x) =
1
2
2
kx (Method B in the figure).
Elastic energy of or within a substance is static energy of configuration. It corresponds to energy stored principally by changing the
inter-atomic distances between nuclei. Thermal energy is the randomized distribution of kinetic energy within the material,
resulting in statistical fluctuations of the material about the equilibrium configuration. There is some interaction, however. For
example, for some solid objects, twisting, bending, and other distortions may generate thermal energy, causing the material’s
temperature to rise. This energy can also produce macroscopic vibrations sufficiently lacking in randomization to lead to
oscillations that are merely the exchange between (elastic) potential energy within the object and the kinetic energy of motion of
the object as a whole.
Key Points
Mathematically, Hooke’s Law can be written as F = −kx .
Many materials obey this law as long as the load does not exceed the material’s elastic limit.
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The rate or spring constant, k, relates the force to the extension in SI units: N/m or kg/s2.
In order to produce a deformation, work must be done.
The potential energy stored in a spring is given by PE = kx , where k is the spring constant and x is the displacement.
el
1
2
2
Deformation can also be converted into thermal energy or cause an object to begin oscillating.
Key Terms
elasticity: The property by virtue of which a material deformed under the load can regain its original dimensions when
unloaded
deformation: A transformation; change of shape.
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
oscillating: Moving in a repeated back-and-forth motion.
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15.3: Periodic Motion
learning objectives
Practice converting between frequency and period
Sinusoidal Waves of Varying Frequencies: Sinusoidal waves of various frequencies; the bottom waves have higher frequencies
than those above. The horizontal axis represents time.
1
f = (15.3.1)
T
For example, if a newborn baby’s heart beats at a frequency of 120 times a minute, its period (the interval between beats) is half a
second. If you calibrate your intuition so that you expect large frequencies to be paired with short periods, and vice versa, you may
avoid some embarrassing mistakes on physics exams.
Units
Locomotive Wheels: The locomotive’s wheels spin at a frequency of f cycles per second, which can also be described as ω radians
per second. The mechanical linkages allow the linear vibration of the steam engine’s pistons, at frequency f, to drive the wheels.
In SI units, the unit of frequency is the hertz (Hz), named after the German physicist Heinrich Hertz: 1 Hz indicates that an event
repeats once per second. A traditional unit of measure used with rotating mechanical devices is revolutions per minute, abbreviated
RPM. 60 RPM equals one hertz (i.e., one revolution per second, or a period of one second). The SI unit for period is the second.
Angular Frequency
Often periodic motion is best expressed in terms of angular frequency, represented by the Greek letter ω (omega). Angular
frequency refers to the angular displacement per unit time (e.g., in rotation) or the rate of change of the phase of a sinusoidal
waveform (e.g., in oscillations and waves), or as the rate of change of the argument of the sine function.
ω = 2πf (15.3.3)
Angular frequency is often represented in units of radians per second (recall there are 2π radians in a circle).
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Period of a Mass on a Spring
−
−
The period of a mass m on a spring of spring constant k can be calculated as T = 2π √ m
k
.
learning objectives
Identify parameters necessary to calculate the period and frequency of an oscillating mass on the end of an ideal spring
Oscillating Ruler: When displaced from its vertical equilibrium position, this plastic ruler oscillates back and forth because of the
restoring force opposing displacement. When the ruler is on the left, there is a force to the right, and vice versa.
Consider, for example, plucking a plastic ruler shown in the first figure. The deformation of the ruler creates a force in the opposite
direction, known as a restoring force. Once released, the restoring force causes the ruler to move back toward its stable equilibrium
position, where the net force on it is zero. However, by the time the ruler gets there, it gains momentum and continues to move to
the right, producing the opposite deformation. It is then forced to the left, back through equilibrium, and the process is repeated
until dissipative forces (e.g., friction) dampen the motion. These forces remove mechanical energy from the system, gradually
reducing the motion until the ruler comes to rest.
Restoring force, momentum, and equilibrium: (a) The plastic ruler has been released, and the restoring force is returning the
ruler to its equilibrium position. (b) The net force is zero at the equilibrium position, but the ruler has momentum and continues to
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move to the right. (c) The restoring force is in the opposite direction. It stops the ruler and moves it back toward equilibrium again.
(d) Now the ruler has momentum to the left. (e) In the absence of damping (caused by frictional forces), the ruler reaches its
original position. From there, the motion will repeat itself.
Hooke’s Law
The simplest oscillations occur when the restoring force is directly proportional to displacement. The name that was given to this
relationship between force and displacement is Hooke’s law:
F = kx (15.3.4)
Here, F is the restoring force, x is the displacement from equilibrium or deformation, and k is a constant related to the difficulty in
deforming the system (often called the spring constant or force constant). Remember that the minus sign indicates the restoring
force is in the direction opposite to the displacement. The force constant k is related to the rigidity (or stiffness) of a system—the
larger the force constant, the greater the restoring force, and the stiffer the system. The units of k are newtons per meter (N/m). For
example, k is directly related to Young’s modulus when we stretch a string. A typical physics laboratory exercise is to measure
restoring forces created by springs, determine if they follow Hooke’s law, and calculate their force constants if they do.
Mass on a Spring
A common example of an objecting oscillating back and forth according to a restoring force directly proportional to the
displacement from equilibrium (i.e., following Hooke’s Law) is the case of a mass on the end of an ideal spring, where “ideal”
means that no messy real-world variables interfere with the imagined outcome.
The motion of a mass on a spring can be described as Simple Harmonic Motion (SHM), the name given to oscillatory motion for a
system where the net force can be described by Hooke’s law. We can now determine how to calculate the period and frequency of
an oscillating mass on the end of an ideal spring. The period T can be calculated knowing only the mass, m, and the force constant,
k:
−
−−
m
T = 2π √ (15.3.5)
k
T
, the frequency is given by:
−
−−
1 k
f = √ (15.3.6)
2π m
We can understand the dependence of these equations on m and k intuitively. If one were to increase the mass on an oscillating
spring system with a given k, the increased mass will provide more inertia, causing the acceleration due to the restoring force F to
decrease (recall Newton’s Second Law: F = ma ). This will lengthen the oscillation period and decrease the frequency. In contrast,
increasing the force constant k will increase the restoring force according to Hooke’s Law, in turn causing the acceleration at each
displacement point to also increase. This reduces the period and increases the frequency. The maximum displacement from
equilibrium is known as the amplitude X.
Motion of a mass on an ideal spring: An object attached to a spring sliding on a frictionless surface is an uncomplicated simple
harmonic oscillator. When displaced from equilibrium, the object performs simple harmonic motion that has an amplitude X and a
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period T. The object’s maximum speed occurs as it passes through equilibrium. The stiffer the spring is, the smaller the period T.
The greater the mass of the object is, the greater the period T. (a) The mass has achieved its greatest displacement X to the right and
now the restoring force to the left is at its maximum magnitude. (b) The restoring force has moved the mass back to its equilibrium
point and is now equal to zero, but the leftward velocity is at its maximum. (c) The mass’s momentum has carried it to its
maximum displacement to the right. The restoring force is now to the right, equal in magnitude and opposite in direction compared
to (a). (d) The equilibrium point is reach again, this time with momentum to the right. (e) The cycle repeats.
learning objectives
Relate the restoring force and the displacement during the simple harmonic motion
Simple Harmonic Motion: A brief introduction to simple harmonic motion for calculus-based physics students.
Simple harmonic motion is typified by the motion of a mass on a spring when it is subject to the linear elastic restoring force given
by Hooke’s Law. A system that follows simple harmonic motion is known as a simple harmonic oscillator.
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Dynamics of Simple Harmonic Oscillation
For one-dimensional simple harmonic motion, the equation of motion (which is a second-order linear ordinary differential equation
with constant coefficients) can be obtained by means of Newton’s second law and Hooke’s law.
2
d x
Fnet = m = −kx, (15.3.7)
2
dt
where m is the mass of the oscillating body, x is its displacement from the equilibrium position, and k is the spring constant.
Therefore:
2
d x k
= −( )x. (15.3.8)
2
dt m
Solving the differential equation above, a solution which is a sinusoidal function is obtained.
x(t) = c1 cos(ωt) + c2 sin(ωt) = A cos(ωt − φ), (15.3.9)
where
−
−−
k
ω =√ , (15.3.10)
m
−−−−−−
2 2
A = √c + c , (15.3.11)
1 2
c2
tan φ = ( ). (15.3.12)
c1
In the solution, c1 and c2 are two constants determined by the initial conditions, and the origin is set to be the equilibrium position.
Each of these constants carries a physical meaning of the motion: A is the amplitude (maximum displacement from the equilibrium
position), ω = 2πf is the angular frequency, and φ is the phase.
We can use differential calculus and find the velocity and acceleration as a function of time:
dx
v(t) = = −Aω sin(ωt − φ) (15.3.13)
dt
2
d x
2
a(t) = = −Aω cos(ωt − φ). (15.3.14)
2
dt
Recalling that T = 1
f
,
−
−−
m
T = 2π √ . (15.3.17)
k
Using Newton’s Second Law, Hooke’s Law, and some differential Calculus, we were able to derive the period and frequency of the
mass oscillating on a spring that we encountered in the last section! Note that the period and frequency are completely independent
of the amplitude.
The below figure shows the simple harmonic motion of an object on a spring and presents graphs of x(t),v(t), and a(t) versus time.
You should learn to create mental connections between the above equations, the different positions of the object on a spring in the
cartoon, and the associated positions in the graphs of x(t), v(t), and a(t).
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Visualizing Simple Harmonic Motion: Graphs of x(t),v(t), and a(t) versus t for the motion of an object on a spring. The net force
on the object can be described by Hooke’s law, and so the object undergoes simple harmonic motion. Note that the initial position
has the vertical displacement at its maximum value X; v is initially zero and then negative as the object moves down; and the initial
acceleration is negative, back toward the equilibrium position and becoming zero at that point.
learning objectives
Describe relationship between the simple harmonic motion and uniform circular motion
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Velocity
The above figure illustrates velocity and acceleration vectors for uniform motion at four different points in the orbit. Since velocity
v is tangent to the circular path, no two velocities point in the same direction. Although the object has a constant speed, its direction
is always changing. This change in velocity is due to an acceleration, a, whose magnitude is (like that of the velocity) held
constant, but whose direction also is always changing. The acceleration points radially inwards (centripetally) and is perpendicular
to the velocity. This acceleration is known as centripetal acceleration.
Uniform Circular Motion (at Four Different Point in the Orbit): Velocity v and acceleration a in uniform circular motion at
angular rate ω; the speed is constant, but the velocity is always tangent to the orbit; the acceleration has constant magnitude, but
always points toward the center of rotation
Displacement around a circular path is often given in terms of an angle θ. This angle is the angle between a straight line drawn
from the center of the circle to the objects starting position on the edge and a straight line drawn from the objects ending position
on the edge to center of the circle. See for a visual representation of the angle where the point p started on the x- axis and moved to
its present position. The angle θ describes how far it moved.
Projection of Uniform Circular Motion: A point P moving on a circular path with a constant angular velocity ω is undergoing
uniform circular motion. Its projection on the x-axis undergoes simple harmonic motion. Also shown is the velocity of this point
around the circle, v−max, and its projection, which is v. Note that these velocities form a similar triangle to the displacement
triangle.
For a path around a circle of radius r, when an angle θ(measured in radians ) is swept out, the distance traveled on the edge of the
circle is s = rθ. You can prove this yourself by remembering that the circumference of a circle is 2πr, so if the object traveled
around the whole circle (one circumference) it will have gone through an angle of 2π radians and traveled a distance of 2πr.
Therefore, the speed of travel around the orbit is:
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dθ
v =r = rω, (15.3.18)
dt
r
. ) Thus,v is a constant, and the velocity vector v also rotates with constant
magnitude v, at the same angular rate ω.
Acceleration
The acceleration in uniform circular motion is always directed inward and is given by:
2
dθ v
a =v = vω = . (15.3.19)
dt r
This acceleration acts to change the direction of v, but not the speed.
Shadow of a Ball Undergoing Simple Harmonic Motion: The shadow of a ball rotating at constant angular velocity ω on a
turntable goes back and forth in precise simple harmonic motion.
The next figure shows the basic relationship between uniform circular motion and simple harmonic motion. The point P travels
around the circle at constant angular velocity ω. The point P is analogous to the ball on a turntable in the figure above. The
projection of the position of P onto a fixed axis undergoes simple harmonic motion and is analogous to the shadow of the object. At
a point in time assumed in the figure, the projection has position x and moves to the left with velocity v. The velocity of the point P
around the circle equals |vmax|. The projection of |vmax| on the x-axis is the velocity v of the simple harmonic motion along the x-
axis.
To see that the projection undergoes simple harmonic motion, note that its position x is given by:
x = X cos θ, (15.3.20)
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where θ=ωt, ω is the constant angular velocity, and X is the radius of the circular path. Thus,
x = X cos ωt. (15.3.21)
The angular velocity ω is in radians per unit time; in this case 2π radians is the time for one revolution T. That is, ω =
2π
T
.
Substituting this expression for ω, we see that the position x is given by:
2πt
x(t) = cos( ) = cos(2πf t). (15.3.22)
T
Note: This equation should look familiar from our earlier discussion of simple harmonic motion.
learning objectives
Identify parameters that affect the period of a simple pendulum
Simple Pendulum: A simple pendulum has a small-diameter bob and a string that has a very small mass but is strong enough not
to stretch appreciably. The linear displacement from equilibrium is s, the length of the arc. Also shown are the forces on the bob,
which result in a net force of −mgsinθ toward the equilibrium position—that is, a restoring force.
For small displacements, a pendulum is a simple harmonic oscillator. A simple pendulum is defined to have an object that has a
small mass, also known as the pendulum bob, which is suspended from a wire or string of negligible mass, such as shown in the
illustrating figure. Exploring the simple pendulum a bit further, we can discover the conditions under which it performs simple
harmonic motion, and we can derive an interesting expression for its period.
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AP Physics C - Pendulums
Pendulums: A brief introduction to pendulums (both ideal and physical) for calculus-based physics students from the standpoint of
simple harmonic motion.
We begin by defining the displacement to be the arc length s. We see from the figure that the net force on the bob is tangent to the
arc and equals −mgsinθ. (The weight mg has components mgcosθ along the string and mgsinθ tangent to the arc. ) Tension in the
string exactly cancels the component mgcosθ parallel to the string. This leaves a net restoring force drawing the pendulum back
toward the equilibrium position at θ = 0.
Now, if we can show that the restoring force is directly proportional to the displacement, then we have a simple harmonic oscillator.
In trying to determine if we have a simple harmonic oscillator, we should note that for small angles (less than about 15º), sinθ≈θ
(sinθ and θ differ by about 1% or less at smaller angles). Thus, for angles less than about 15º, the restoring force F is
F ≈ −mgθ. (15.3.23)
The displacement s is directly proportional to θ. When θ is expressed in radians, the arc length in a circle is related to its radius (L
in this instance) by:
s=Lθs=Lθ
so that
θ = sL. (15.3.24)
For small angles, then, the expression for the restoring force is:
mgL
F ≈ . (15.3.25)
s
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This expression is of the form of Hooke’s Law:
F ≈ −kx (15.3.26)
where the force constant is given by k=mg/L and the displacement is given by x=s. For angles less than about 15º, the restoring
force is directly proportional to the displacement, and the simple pendulum is a simple harmonic oscillator.
Using this equation, we can find the period of a pendulum for amplitudes less than about 15º. For the simple pendulum:
−
−− −−−−
m m
T = 2π √ = 2π . (15.3.27)
mg
k
⎷
L
Thus,
−−
L
T = 2π √ (15.3.28)
g
or the period of a simple pendulum. This result is interesting because of its simplicity. The only things that affect the period of a
simple pendulum are its length and the acceleration due to gravity. The period is completely independent of other factors, such as
mass. Even simple pendulum clocks can be finely adjusted and accurate. Note the dependence of T on g. If the length of a
pendulum is precisely known, it can actually be used to measure the acceleration due to gravity. If θ is less than about 15º, the
period T for a pendulum is nearly independent of amplitude, as with simple harmonic oscillators. In this case, the motion of a
pendulum as a function of time can be modeled as:
2πt
θ(t) = θo cos( ) (15.3.29)
T
For amplitudes larger than 15º, the period increases gradually with amplitude so it is longer than given by the simple equation for T
above. For example, at an amplitude of θ0 = 23° it is 1% larger. The period increases asymptotically (to infinity) as θ0approaches
180°, because the value θ0 = 180° is an unstable equilibrium point for the pendulum.
learning objectives
Identify parameters that affect the period of a physical pendulum
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AP Physics C - Pendulums
Pendulums – Physical Pendulum: A brief introduction to pendulums (both ideal and physical) for calculus-based physics students
from the standpoint of simple harmonic motion.
A Physical Pendulum: An example showing how forces act through center of mass. We can calculate the period of this pendulum
by determining the moment of inertia of the object around the pivot point.
Gravity acts through the center of mass of the rigid body. Hence, the length of the pendulum used in equations is equal to the linear
distance between the pivot and the center of mass (h).
The equation of torque gives:
τ = Iα, (15.3.30)
where α is the angular acceleration, τ is the torque, and I is the moment of inertia.
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The torque is generated by gravity so:
τ = mgh sin θ, (15.3.31)
where h is the distance from the center of mass to the pivot point and θ is the angle from the vertical.
Hence, under the small-angle approximation sin\theta \approx \theta,
mghθ
α ≈− . (15.3.32)
I
This is of the same form as the conventional simple pendulum and this gives a period of:
−−−−
I
T = 2π √ . (15.3.33)
mgh
Uniform Rigid Rod: A rigid rod with uniform mass distribution hangs from a pivot point. This is another example of a physical
pendulum.
L
h= . (15.3.34)
2
The moment of inertia of the rigid rod about its center is:
2
mL
Ic = . (15.3.35)
12
However, we need to evaluate the moment of inertia about the pivot point, not the center of mass, so we apply the parallel axis
theorem:
2 2
2
mL L 2
mL
Io = Ic + mh = + m( ) = . (15.3.36)
12 2 3
The important thing to note about this relation is that the period is still independent of the mass of the rigid body. However, it is not
independent of the mass distribution of the rigid body. A change in shape, size, or mass distribution will change the moment of
inertia. This, in turn, will change the period.
As with a simple pendulum, a physical pendulum can be used to measure g.
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Energy in a Simple Harmonic Oscillator
The total energy in a simple harmonic oscillator is the constant sum of the potential and kinetic energies.
learning objectives
Explain why the total energy of the harmonic oscillator is constant
where PE is the potential energy, k is the spring constant, and x is the magnitude of the displacement or deformation. Because a
simple harmonic oscillator has no dissipative forces , the other important form of energy is kinetic energy (KE). Conservation of
energy for these two forms is:
KE + PE = constant, (15.3.39)
This statement of conservation of energy is valid for all simple harmonic oscillators, including ones where the gravitational force
plays a role. For example, for a simple pendulum we replace the velocity with v=Lω, the spring constant with k=mg/L, and the
displacement term with x=Lθ. Thus:
1 2 2
1 2
mL ω + mgLθ = constant. (15.3.41)
2 2
In the case of undamped, simple harmonic motion, the energy oscillates back and forth between kinetic and potential, going
completely from one to the other as the system oscillates. So for the simple example of an object on a frictionless surface attached
to a spring, as shown again (see ), the motion starts with all of the energy stored in the spring. As the object starts to move, the
elastic potential energy is converted to kinetic energy, becoming entirely kinetic energy at the equilibrium position. It is then
converted back into elastic potential energy by the spring, the velocity becomes zero when the kinetic energy is completely
converted, and so on. This concept provides extra insight here and in later applications of simple harmonic motion, such as
alternating current circuits.
Energy in a Simple Harmonic Oscillator: The transformation of energy in simple harmonic motion is illustrated for an object
attached to a spring on a frictionless surface. (a) The mass has achieved maximum displacement from equilibrium. All energy is
potential energy. (b) As the mass passes through the equilibrium point with maximum speed all energy in the system is in kinetic
energy. (c) Once again, all energy is in the potential form, stored in the compression of the spring (in the first panel the energy was
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stored in the extension of the spring). (d) Passing through equilibrium again all energy is kinetic. (e) The mass has completed an
entire cycle.
The conservation of energy principle can be used to derive an expression for velocity v. If we start our simple harmonic motion
with zero velocity and maximum displacement (x=X), then the total energy is:
1 2
E = kX . (15.3.42)
2
This total energy is constant and is shifted back and forth between kinetic energy and potential energy, at most times being shared
by each. The conservation of energy for this system in equation form is thus:
1 2
1 2
1 2
mv + kx = kX . (15.3.43)
2 2 2
and so:
−−−−−−
2
x
v = ±vmax √ 1 − , (15.3.46)
2
X
where:
−
−−
k
vmax = √ X. (15.3.47)
m
From this expression, we see that the velocity is a maximum (vmax) at x=0. Notice that the maximum velocity depends on three
factors. It is directly proportional to amplitude. As you might guess, the greater the maximum displacement, the greater the
maximum velocity. It is also greater for stiffer systems because they exert greater force for the same displacement. This observation
is seen in the expression for vmax; it is proportional to the square root of the force constant k. Finally, the maximum velocity is
smaller for objects that have larger masses, because the maximum velocity is inversely proportional to the square root of m. For a
given force, objects that have large masses accelerate more slowly.
A similar calculation for the simple pendulum produces a similar result, namely:
−−
g
ωmax = √ θmax . (15.3.48)
L
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Joy Physics- exp #13 Spring-mass osc
Experience with a simple harmonic oscillator: A known mass is hung from a spring of known spring constant and allowed to
oscillate. The time for one oscillation (period) is measured. This value is compared to a predicted value, based on the mass and
spring constant.
learning objectives
Review factors responsible for the sinusoidal behavior of uniform circular motion
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Sinusoidal and Non-Sinusoidal Vibrations: Only the top graph is sinusoidal. The others vary with constant amplitude and period,
but do no describe simple harmonic motion.
When this general equation is solved for the position, velocity and acceleration as a function of time:
x(t) = A cos(ωt − φ)
dx
v(t) = = −Aω sin(ωt − φ)
dt
2
d x 2
a(t) = 2
= −Aω cos(ωt − φ)
dt
These are all sinusoidal solutions. Consider a mass on a spring that has a small pen inside running across a moving strip of paper as
it bounces, recording its movements.
Mass on Spring Producing Sine Wave: The vertical position of an object bouncing on a spring is recorded on a strip of moving
paper, leaving a sine wave.
The above equations can be rewritten in a form applicable to the variables for the mass on spring system in the figure.
2πt
x(t) = X cos( )
T
2πt
v(t) = −vmax sin( )
T
kX 2πt
a(t) = − cos( )
m T
Recall that the projection of uniform circular motion can be described in terms of a simple harmonic oscillator. Uniform circular
motion is therefore also sinusoidal, as you can see from.
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Sinusoidal Nature of Uniform Circular Motion: The position of the projection of uniform circular motion performs simple
harmonic motion, as this wavelike graph of x versus t indicates.
Summing K(t) and U(t) produces the total mechanical energy seen before:
1
2
E = K+U = kA . (15.3.52)
2
Key Points
Motion that repeats itself regularly is called periodic motion. One complete repetition of the motion is called a cycle. The
duration of each cycle is the period.
The frequency refers to the number of cycles completed in an interval of time. It is the reciprocal of the period and can be
calculated with the equation f=1/T.
Some motion is best characterized by the angular frequency (ω). The angular frequency refers to the angular displacement per
unit time and is calculated from the frequency with the equation ω = 2πf .
If an object is vibrating to the right and left, then it must have a leftward force on it when it is on the right side, and a rightward
force when it is on the left side.
The restoring force causes an oscillating object to move back toward its stable equilibrium position, where the net force on it is
zero.
The simplest oscillations occur when the restoring force is directly proportional to displacement. In this case the force can be
calculated as F = −kx , where F is the restoring force, k is the force constant, and x is the displacement.
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The motion of a mass on a spring can be described as Simple Harmonic Motion (SHM): oscillatory motion that follows Hooke’s
Law.
−
−
The period of a mass on a spring is given by the equation T = 2π √ m
Simple harmonic motion is often modeled with the example of a mass on a spring, where the restoring force obey’s Hooke’s
Law and is directly proportional to the displacement of an object from its equilibrium position.
Any system that obeys simple harmonic motion is known as a simple harmonic oscillator.
The equation of motion that describes simple harmonic motion can be obtained by combining Newton’s Second Law and
2
Uniform circular motion describes the movement of an object traveling a circular path with constant speed. The one-
dimensional projection of this motion can be described as simple harmonic motion.
In uniform circular motion, the velocity vector v is always tangent to the circular path and constant in magnitude. The
acceleration is constant in magnitude and points to the center of the circular path, perpendicular to the velocity vector at every
instant.
If an object moves with angular velocity ω around a circle of radius r centered at the origin of the x-y plane, then its motion
along each coordinate is simple harmonic motion with amplitude r and angular frequency ω.
A simple pendulum is defined as an object that has a small mass, also known as the pendulum bob, which is suspended from a
wire or string of negligible mass.
When displaced, a pendulum will oscillate around its equilibrium point due to momentum in balance with the restoring force of
gravity.
When the swings ( amplitudes ) are small, less than about 15º, the pendulum acts as a simple harmonic oscillator with period
−
−
T = 2π √
L
g
, where L is the length of the string and g is the acceleration due to gravity.
A physical pendulum is the generalized case of the simple pendulum. It consists of any rigid body that oscillates about a pivot
point.
For small amplitudes, the period of a physical pendulum only depends on the moment of inertia of the body around the pivot
−−−
point and the distance from the pivot to the body’s center of mass. It is calculated as:T = 2π √ mgh
I
.
The period is still independent of the total mass of the rigid body. However, it is not independent of the mass distribution of the
rigid body. A change in shape, size, or mass distribution will change the moment of inertia and thus, the period.
The sum of the kinetic and potential energies in a simple harmonic oscillator is a constant, i.e., KE + PE = constant . The
energy oscillates back and forth between kinetic and potential, going completely from one to the other as the system oscillates.
In a spring system, the conservation equation is written as: mv + kx = constant = kX , where X is the maximum
1
2
2 1
2
2 1
2
2
displacement.
−
−
The maximum velocity depends on three factors: amplitude, the stiffness factor, and mass: v max =√
k
m
X .
For simple harmonic oscillators, the equation of motion is always a second order differential equation that relates the
acceleration and the displacement. The relevant variables are x, the displacement, and k, the spring constant.
Solving the differential equation above always produces solutions that are sinusoidal in nature. For example,
x(t), v(t), a(t), K(t), and U(t) all have sinusoidal solutions for simple harmonic motion.
Uniform circular motion is also sinusoidal because the projection of this motion behaves like a simple harmonic oscillator.
Key Terms
period: The duration of one cycle in a repeating event.
angular frequency: The angular displacement per unit time.
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time t in which it occurs: f = n / t.
Restoring force: A variable force that gives rise to an equilibrium in a physical system. If the system is perturbed away from
the equilibrium, the restoring force will tend to bring the system back toward equilibrium. The restoring force is a function only
of position of the mass or particle. It is always directed back toward the equilibrium position of the system
amplitude: The maximum absolute value of some quantity that varies.
simple harmonic oscillator: A device that implements Hooke’s law, such as a mass that is attached to a spring, with the other
end of the spring being connected to a rigid support, such as a wall.
oscillator: A pattern that returns to its original state, in the same orientation and position, after a finite number of generations.
15.3.19 https://phys.libretexts.org/@go/page/14531
centripetal acceleration: Acceleration that makes a body follow a curved path: it is always perpendicular to the velocity of a
body and directed towards the center of curvature of the path.
uniform circular motion: Movement around a circular path with constant speed.
simple pendulum: A hypothetical pendulum consisting of a weight suspended by a weightless string.
physical pendulum: A pendulum where the rod or string is not massless, and may have extended size; that is, an arbitrarily-
shaped, rigid body swinging by a pivot. In this case, the pendulum’s period depends on its moment of inertia around the pivot
point.
mass distribution: Describes the spatial distribution, and defines the center, of mass in an object.
elastic potential energy: The energy stored in a deformable object, such as a spring.
dissipative forces: Forces that cause energy to be lost in a system undergoing motion.
sinusoidal: In the form of a wave, especially one whose amplitude varies in proportion to the sine of some variable (such as
time).
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Uniform circular motion. Provided by: Wikipedia. Located at: http://en.Wikipedia.org/wiki/Uniform_circular_motion.
License: Public Domain: No Known Copyright
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Pendulums. Located at: http://www.youtube.com/watch?v=gk4KrcKIQ50. License: Public Domain: No Known Copyright.
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Experience with a simple harmonic oscillator. Located at: http://www.youtube.com/watch?v=Iuv24zcc5kI. License: Public
Domain: No Known Copyright. License Terms: Standard YouTube license
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15.4: Damped and Driven Oscillations
learning objectives
Describe the time evolution of the motion of the damped harmonic oscillator
Damped Harmonic Motion: Illustrating the position against time of our object moving in simple harmonic motion. We see that for
small damping, the amplitude of our motion slowly decreases over time.
The simplest and most commonly seen case occurs when the frictional force is proportional to an object’s velocity. Note that other
cases exist which may lead to nonlinear equations which go beyond the scope of this example.
Consider an object of mass m attached to a spring of constant k. Let the damping force be proportional to the mass’ velocity by a
proportionality constant, b, called the vicious damping coefficient. We can describe this situation using Newton’s second law, which
leads to a second order, linear, homogeneous, ordinary differential equation. We simply add a term describing the damping force to
our already familiar equation describing a simple harmonic oscillator to describe the general case of damped harmonic motion.
2
d x dx
Fnet = m +b + kx = 0 (15.4.1)
2
dt dt
2
d x b dx k
= + + x =0 (15.4.2)
2
dt m dt m
2
d x dx
2
= +γ +ω x = 0 (15.4.3)
2 0
dt dt
k b
2
ω = ,γ = (15.4.4)
0
m m
dt
2
dx
dt
, the velocity of our object, ω , undamped angular frequency of oscillation,
0
Plugging this into the differential equation we find that there are three results for a, which will dictate the motion of our system. We
can solve for a by using the quadratic equation.
2 2
Fnet = a x + γax + ω x = 0 (15.4.6)
0
2 2
=a + γa + ω =0 (15.4.7)
0
−−−−−−−
2
γ ± √ γ2 − 4 ω
0
a = (15.4.8)
2
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The physical situation has three possible results depending on the value of a, which depends on the value of what is under our
radical. This expression can be positive, negative, or equal to zero which will result in overdamping, underdamping, and critical
damping, respectively.
γ
2
> 4ω
2
0
is the Over Damped case. In this case, the system returns to equilibrium by exponentially decaying towards zero. The
system will not pass the equilibrium position more than once.
γ
2
< 4ω
2
0
is the Under Damped case. In this case, the system oscillates as it slowly returns to equilibrium and the amplitude
decreases over time. Figure 1 depicts an underdamped case.
γ
2
= 4ω
2
0
is theCritically Damped case. In this case, the system returns to equilibrium very quickly without oscillating and without
passing the equilibrium position at all.
learning objectives
Describe a driven harmonic oscillator as a type of damped oscillator
In classical mechanics, a harmonic oscillator is a system that, when displaced from its equilibrium position, experiences a restoring
→
force, F, proportional to the displacement, xF →= −k x where k is a positive constant. If a frictional force ( damping )
proportional to the velocity is also present, the harmonic oscillator is described as a damped oscillator.
Driven harmonic oscillators are damped oscillators further affected by an externally applied force F(t). Newton’s second law takes
2 2
F(t)
the form F(t) − kx − c
dx
dt
=m
d x
dt
2
. It is usually rewritten into the form d x
dt
2
+ 2ζω0
dx
dt
2
+ω x =
0 m
. This equation can be
2
solved exactly for any driving force, using the solutions z(t) which satisfy the unforced equation: d z
dt
2
+ 2ζω0
dz
dt
2
+ω z = 0
0
, and
−−−−−
which can be expressed as damped sinusoidal oscillationsz(t) = Ae −ζω0 t
sin (√1 − ζ
2
ω0 t + φ) in the case where ζ ≤1 . The
amplitude A and phase φ determine the behavior needed to match the initial conditions. In the case ζ < 1 and a unit step input with
2
sin ( √1−ζ ω0 t+φ)
sin (φ)
with phase φ given by cos φ = ζ . The time an oscillator needs to adapt
to changed external conditions is of the order \(\mathrm{τ =\frac{1}{(ζ_0)}}\). In physics, the adaptation is called relaxation, and τ
is called the relaxation time.
2
dt
2
+ω x =
0
1
m
F0 sin (ωt) , where F is the driving amplitude and ωω is the
0
dt
driving frequency for a sinusoidal driving mechanism. This type of system appears in AC driven RLC circuits (resistor-inductor-
capacitor) and driven spring systems having internal mechanical resistance or external air resistance. The general solution is a sum
of a transient solution that depends on initial conditions, and a steady state that is independent of initial conditions and depends
only on the driving amplitude F , driving frequency ωω, undamped angular frequency ω , and the damping ratio ζ . For a particular
0 0
−−−−−−
driving frequency called the resonance, or resonant frequency ωr = ω0 √1 − 2ζ
2
, the amplitude (for a given F0 ) is maximum.
–
This resonance effect only occurs when , i.e. for significantly underdamped systems. For strongly underdamped systems
ζ < 1/ √2
the value of the amplitude can become quite large near the resonance frequency (see ).
15.4.2 https://phys.libretexts.org/@go/page/14532
Resonance: Steady state variation of amplitude with frequency and damping of a driven simple harmonic oscillator.
Key Points
To describe a damped harmonic oscillator, add a velocity dependent term, bx, where b is the vicious dampingcoefficient.
Solve the differential equation for the equation of motion, x(t).
Depending on the values of the damping coefficient and undamped angular frequency, the results will be one of three cases: an
under damped system, an over damped system, or a critically damped system.
2
dt
d x
for driven harmonic oscillators.
2
dt
Key Terms
Under Damped: “The condition in which damping of an oscillator causes it to return to equilibrium with the amplitude
gradually decreasing to zero; system returns to equilibrium faster but overshoots and crosses the equilibrium position one or
more times. “
Critically Damped: “The condition in which the damping of an oscillator causes it to return as quickly as possible to its
equilibrium position without oscillating back and forth about this position. “
Over Damped: “The condition in which damping of an oscillator causes it to return to equilibrium without oscillating;
oscillator moves more slowly toward equilibrium than in the critically damped system. “
oscillator: A pattern that returns to its original state, in the same orientation and position, after a finite number of generations.
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
force: A physical quantity that denotes ability to push, pull, twist or accelerate a body which is measured in a unit dimensioned
in mass × distance/time² (ML/T²): SI: newton (N); CGS: dyne (dyn)
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15.5: Waves
learning objectives
Describe process of energy and mass transfer during wave motion
Vibrations and waves are extremely important phenomena in physics. In nature, oscillations are found everywhere. From the
jiggling of atoms to the large oscillations of sea waves, we find examples of vibrations in almost every physical system. In physics
a wave can be thought of as a disturbance or oscillation that travels through space-time, accompanied by a transfer of energy. Wave
motion transfers energy from one point to another, often with no permanent displacement of the particles of the medium —that is,
with little or no associated mass transport. They consist, instead, of oscillations or vibrations around almost fixed locations.
The emphasis of the last point highlights an important misconception of waves. Waves transfer energy not mass. An easy way to
see this is to imagine a floating ball a few yards out to sea. As the waves propagate (i.e., travel) towards the shore, the ball will not
come towards the shore. It may come to shore eventually due to the tides, current or wind, but the waves themselves will not carry
the ball with them. A wave only moves mass perpendicular to the direction of propagation—in this case up and down, as illustrated
in the figure below:
Wave motion: The point along the axis is analogous to the floating ball at sea. We notice that while it moves up and down it does
not move in the direction of the wave’s propagation.
A wave can be transverse or longitudinal depending on the direction of its oscillation. Transverse waves occur when a disturbance
causes oscillations perpendicular (at right angles) to the propagation (the direction of energy transfer). Longitudinal waves occur
when the oscillations are parallel to the direction of propagation. While mechanical waves can be both transverse and longitudinal,
all electromagnetic waves are transverse. Sound, for example, is a longitudinal wave.
The description of waves is closely related to their physical origin for each specific instance of a wave process. For example,
acoustics is distinguished from optics in that sound waves are related to a mechanical rather than an electromagnetic (light) wave
transfer caused by vibration. Therefore, concepts such as mass, momentum, inertia or elasticity become crucial in describing
acoustic (as distinct from optic) wave processes. This difference in origin introduces certain wave characteristics particular to the
properties of the medium involved. In this chapter we will closely examine the difference between longitudinal and transverse
waves along with some of the properties they possess. We will also learn how waves are fundamental in describing motion of many
applicable physical systems.
15.5.1 https://phys.libretexts.org/@go/page/14533
High School Physics - The Wave Equation
The Wave Equation: A brief introduction to the wave equation, discussing wave velocity, frequency, wavelength, and period.
Transverse Waves
Transverse waves propagate through media with a speed →vwv→w orthogonally to the direction of energy transfer.
learning objectives
Describe properties of the transverse wave
A transverse wave is a moving wave that consists of oscillations occurring perpendicular (or right angled) to the direction of energy
transfer. If a transverse wave is moving in the positive x-direction, its oscillations are in up and down directions that lie in the y–z
plane. Light is an example of a transverse wave. For transverse waves in matter, the displacement of the medium is perpendicular
to the direction of propagation of the wave. A ripple on a pond and a wave on a string are easily visualized transverse waves.
Transverse waves are waves that are oscillating perpendicularly to the direction of propagation. If you anchor one end of a ribbon
or string and hold the other end in your hand, you can create transverse waves by moving your hand up and down. Notice though,
that you can also launch waves by moving your hand side-to-side. This is an important point. There are two independent directions
in which wave motion can occur. In this case, these are the y and z directions mentioned above. depicts the motion of a transverse
15.5.2 https://phys.libretexts.org/@go/page/14533
wave. Here we observe that the wave is moving in t and oscillating in the x-y plane. A wave can be thought as comprising many
particles (as seen in the figure) which oscillate up and down. In the figure we observe this motion to be in x-y plane (denoted by the
red line in the figure). As time passes the oscillations are separated by units of time. The result of this separation is the sine curve
we expect when we plot position versus time.
Sine Wave: The direction of propagation of this wave is along the t axis.
When a wave travels through a medium–i.e., air, water, etc., or the standard reference medium (vacuum)–it does so at a given
speed: this is called the speed of propagation. The speed at which the wave propagates is denoted and can be found using the
following formula:
v = fλ (15.5.1)
where v is the speed of the wave, f is the frequency , and is the wavelength. The wavelength spans crest to crest while the amplitude
is 1/2 the total distance from crest to trough. Transverse waves have their applications in many areas of physics. Examples of
transverse waves include seismic S (secondary) waves, and the motion of the electric (E) and magnetic (M) fields in an
electromagnetic plane waves, which both oscillate perpendicularly to each other as well as to the direction of energy transfer.
Therefore an electromagnetic wave consists of two transverse waves, visible light being an example of an electromagnetic wave.
Wavelength and Amplitude: The wavelength is the distance between adjacent crests. The amplitude is the 1/2 the distance from
crest to trough.
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Longitudinal vs. Transverse | Two Types of Waves | …
Longitudinal Waves
Longitudinal waves, sometimes called compression waves, oscillate in the direction of propagation.
learning objectives
Give properties and provide examples of the longitudinal wave
Longitudinal Waves
Longitudinal waves have the same direction of vibration as their direction of travel. This means that the movement of the medium
is in the same direction as the motion of the wave. Some longitudinal waves are also called compressional waves or compression
waves. An easy experiment for observing longitudinal waves involves taking a Slinky and holding both ends. After compressing
and releasing one end of the Slinky (while still holding onto the end), a pulse of more concentrated coils will travel to the end of the
Slinky.
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Longitudinal Waves: A compressed Slinky is an example of a longitudinal wave. The wve propagates in the same direction of
oscillation.
Like transverse waves, longitudinal waves do not displace mass. The difference is that each particle which makes up the medium
through which a longitudinal wave propagates oscillates along the axis of propagation. In the example of the Slinky, each coil will
oscillate at a point but will not travel the length of the Slinky. It is important to remember that energy, in this case in the form of a
pulse, is being transmitted and not the displaced mass.
Longitudinal waves can sometimes also be conceptualized as pressure waves. The most common pressure wave is the sound wave.
Sound waves are created by the compression of a medium, usually air. Longitudinal sound waves are waves of alternating pressure
deviations from the equilibrium pressure, causing local regions of compression and rarefaction. Matter in the medium is
periodically displaced by a sound wave, and thus oscillates. When people make a sound, whether it is through speaking or hitting
something, they are compressing the air particles to some significant amount. By doing so, they create transverse waves. When
people hear sounds, their ears are sensitive to the pressure differences and interpret the waves as different tones.
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Two Types of Waves: Longitudinal vs. Transverse: Even ocean waves!
Water Waves
Water waves can be commonly observed in daily life, and comprise both transverse and longitudinal wave motion.
learning objectives
Describe particle movement in water waves
Water waves, which can be commonly observed in our daily lives, are of specific interest to physicists. Describing detailed fluid
dynamics in water waves is beyond the scope of introductory physics courses. Although we often observe water wave propagating
in 2D, in this atom we will limit our discussion to 1D propagation.
Plane wave: We see a wave propagating in the direction of the phase velocity. The wave can be thought to be made up of planes
orthogonal to the direction of the phase velocity.
Since water waves transport energy, attempts to generate power from them have been made by utilizing the physical motion of such
waves. Although larger waves are more powerful, wave power is also determined by wave speed, wavelength, and water density.
Deep water corresponds with a water depth larger than half the wavelength, as is a common case in the sea and ocean. In deep
water, longer-period waves propagate faster and transport their energy faster. The deep-water group velocity is half the phase
velocity. In shallow water for wavelengths larger than about twenty times the water depth (as often found near the coast), the group
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velocity is equal to the phase velocity. These methods have proven viable in some cases but do not provide a fully sustainable form
of renewable energy to date.
Water waves: The motion water waves causes particles to follow clockwise circular motion. This is a result of the wave having
both transverse and longitudinal properties.
learning objectives
Identify major characteristic properties of waves
Characteristics of Waves
Waves have certain characteristic properties which are observable at first notice. The first property to note is the amplitude. The
amplitude is half of the distance measured from crest to trough. We also observe the wavelength, which is the spatial period of the
wave (e.g. from crest to crest or trough to trough). We denote the wavelength by the Greek letter λλ.
The frequency of a wave is the number of cycles per unit time — one can think of it as the number of crests which pass a fixed
point per unit time. Mathematically, we make the observation that,
Frequencies of different sine waves.: The red wave has a low frequency sine there is very little repetition of cycles. Conversely
we say that the purple wave has a high frequency. Note that time increases along the horizontal.
1
f = (15.5.2)
T
where T is the period of oscillation. Frequency and wavelength can also be related-* with respects to a “speed” of a wave. In fact,
\[\mathrm{v=fλ}]
where v is called the wave speed, or more commonly,the phase velocity, the rate at which the phase of the wave propagates in
space. This is the velocity at which the phase of any one frequency component of the wave travels. For such a component, any
given phase of the wave (for example, the crest) will appear to travel at the phase velocity.
Finally, the group velocity of a wave is the velocity with which the overall shape of the waves’ amplitudes — known as the
modulation or envelope of the wave — propagates through space. In, one may see that the overall shape (or “envelope”) propagates
to the right, while the phase velocity is negative.
15.5.7 https://phys.libretexts.org/@go/page/14533
Fig 2: This shows a wave with the group velocity and phase velocity going in different directions. (The group velocity is positive
and the phase velocity is negative. )
Energy Transportation
Waves transfer energy which can be used to do work.
learning objectives
Relate direction of energy and wave transportation
Energy transportion is essential to waves. It is a common misconception that waves move mass. Waves carry energy along an axis
defined to be the direction of propagation. One easy example is to imagine that you are standing in the surf and you are hit by a
significantly large wave, and once you are hit you are displaced (unless you hold firmly to your ground!). In this sense the wave
has done work (it applied a force over a distance). Since work is done over time, the energy carried by a wave can be used to
generate power.
Water Wave: Waves that are more massive or have a greater velocity transport more energy.
Similarly we find that electromagnetic waves carry energy. Electromagnetic radiation (EMR) carries energy—sometimes called
radiant energy—through space continuously away from the source (this is not true of the near-field part of the EM field).
Electromagnetic waves can be imagined as a self-propagating transverse oscillating wave of electric and magnetic fields. EMR also
carries both momentum and angular momentum. These properties may all be imparted to matter with which it interacts (through
work). EMR is produced from other types of energy when created, and it is converted to other types of energy when it is destroyed.
The photon is the quantum of the electromagnetic interaction, and is the basic “unit” or constituent of all forms of EMR. The
quantum nature of light becomes more apparent at high frequencies (or high photon energy). Such photons behave more like
particles than lower-frequency photons do.
15.5.8 https://phys.libretexts.org/@go/page/14533
Electromagnetic Wave: Electromagnetic waves can be imagined as a self-propagating transverse oscillating wave of electric and
magnetic fields. This 3D diagram shows a plane linearly polarized wave propagating from left to right.
In general, there is a relation of waves which states that the velocity (v ) of a wave is proportional to the frequency (f ) times the
wavelength (λ ):
v = fλ (15.5.3)
dp
We also know that classical momentum pp is given by p=mvp=mv which relates to force via Newton’s second law: F =
dt
EM waves with higher frequencies carry more energy. This is a direct result of the equations above. Since v ∝ f we find that higher
frequencies imply greater velocity. If velocity is increased then we have greater momentum which implies a greater force (it gets a
little bit tricky when we talk about particles moving close to the speed of light, but this observation holds in the classical sense).
Since energy is the ability of an object to do work, we find that for W = Fd a greater force correlates to more energy transfer.
Again, this is an easy phenomenon to experience empirically; just stand in front of a faster wave and feel the difference!
Key Points
A wave can be thought of as a disturbance or oscillation that travels through space-time, accompanied by a transfer of energy.
The direction a wave propagates is perpendicular to the direction it oscillates for transverse waves.
A wave does not move mass in the direction of propagation; it transfers energy.
Transverse waves oscillate in the z-y plane but travel along the x axis.
A transverse wave has a speed of propagation given by the equation v = f λ .
The direction of energy transfer is perpendicular to the motion of the wave.
While longitudinal waves oscillate in the direction of propagation, they do not displace mass since the oscillations are small and
involve an equilibrium position.
The longitudinal ‘waves’ can be conceptualized as pulses that transfer energy along the axis of propagation.
Longitudinal waves can be conceptualized as pressure waves characterized by compression and rarefaction.
The particles which make up a water wave move in circular paths.
If the waves move slower than the wind above them, energy is transfered from the wind to the waves.
The oscillations are greatest on the surface of the wave and become weaker deeper in the fluid.
The wavelength is the spatial period of the wave.
The frequency of a wave refers to the number of cycles per unit time and is not to be confused with angular frequency.
The phase velocity can be expressed as the product of wavelength and frequency.
Waves which are more massive transfer more energy.
Waves with greater velocities transfer more energy.
Energy of a wave is transported in the direction of the waves transportation.
Key Terms
medium: The material or empty space through which signals, waves or forces pass.
direction of propagation: The axis along which the wave travels.
wave: A moving disturbance in the energy level of a field.
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ, and corresponds to the velocity of the wave divided by its frequency.
trough: A long, narrow depression between waves or ridges.
speed of propagation: The speed at which a wave moves through a medium.
crest: The ridge or top of a wave.
transverse wave: Any wave in which the direction of disturbance is perpendicular to the direction of travel.
rarefaction: a reduction in the density of a material, especially that of a fluid
Longitudinal: Running in the direction of the long axis of a body.
compression: to increase in density; the act of compressing, or the state of being compressed; compaction
phase velocity: The velocity of propagation of a pure sine wave of infinite extent and infinitesimal amplitude.
group velocity: The propagation velocity of the envelope of a modulated travelling wave, which is considered as the
propagation velocity of information or energy contained in it.
15.5.9 https://phys.libretexts.org/@go/page/14533
plane wave: A constant-frequency wave whose wavefronts (surfaces of constant phase) are infinite parallel planes of constant
peak-to-peak amplitude normal to the phase velocity vector.
wave speed: The absolute value of the velocity at which the phase of any one frequency component of the wave travels.
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ, and corresponds to the velocity of the wave divided by its frequency.
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time t in which it occurs: f = . n
energy: A quantity that denotes the ability to do work and is measured in a unit dimensioned in mass × distance²/time² (ML²/T²)
or the equivalent.
power: A measure of the rate of doing work or transferring energy.
work: A measure of energy expended in moving an object; most commonly, force times displacement. No work is done if the
object does not move.
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15.6: Wave Behavior and Interaction
learning objectives
Distinguish transmission and reflection phenomena
When the medium through which a wave travels suddenly changes, the wave often experiences partial transmission and partial
refection at the interface. Reflection is a wave phenomenon that changes the direction of a wavefront at an interface between two
different media so that the wavefront returns into the medium from which it originated. Transmission permits the passage of wave,
with some or none of the incident wave being absorbed. Reflection and transmission often occur at the same time.
Partial Transmittance and Partial Reflectance: A wave experiences partial transmittance and partial reflectance when the
medium through which it travels suddenly changes.
Consider a long string made by connecting two sub-strings with different density μ , μ . When the string is driven by an external
1 2
force, partial reflection and transmission occur as in the figure above. For the incoming, reflected, and transmitted waves, we can
try a solution of the following forms:
Two Strings With Different Density: Two strings with different density are connected and driven by an external driving force.
yinc = A cos(k1 x − ωt) (15.6.1)
k1 and k2 are determined by the speed of the wave in each medium. We choose our coordinates such that the junction of two sub-
strings is located at x=0. In choosing a trial solution for the waves, we assumed that the incident and transmitted waves travel to the
right, while the reflected waves travel to the left. (This is why the ‘+’ sign is chosen before ωtωt in the reflected wave. On the left
side of the junction, we have
yl = yinc + yref = A cos(k1 x − ωt). (15.6.4)
We will impose additional restriction on the waves by applying “boundary conditions” at x=0. At the boundary x=0, the wave must
be continuous and there should be no kinks in it. Thus we must have
yl (x = 0, t) = yr (x = 0, t) (15.6.6)
k2
From the first equation, we get A + B = C . From the second equation, we get A − B = ( k1
)C .
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1 k2
A = (1 + )C (15.6.8)
2 k1
1 k2
B= (1 − )C (15.6.9)
2 k1
learning objectives
Distinguish destructive and constructive interference and identify conditions that are required for the superposition of
waves
Most waves do not look very simple. They look are often more complex than the simple water waves often considered in
textbooks. Simple waves may be created by a simple harmonic oscillation, and thus have a sinusoidal shape. Complex waves are
more interesting, even beautiful, but they look formidable. Most waves appear complex because they result from several simple
waves adding together. Luckily, the rules for adding waves are quite simple.
Complex Pattern of Waves: These waves result from the superposition of several waves from different sources, producing a
complex pattern.
Superposition
When two or more waves arrive at the same point, they superimpose themselves on one another. More specifically, the disturbances
of waves are superimposed when they come together—a phenomenon called superposition. Each disturbance corresponds to a
force, and forces add. If the disturbances are along the same line, then the resulting wave is a simple addition of the disturbances of
the individual waves—that is, their amplitudes add.
Interference
As a result of superposition of waves, interference can be observed. Interference is an effect caused by two or more waves.
15.6.2 https://phys.libretexts.org/@go/page/14534
High School Physics - Wave Interference
Wave Interference: A brief introduction to constructive and destructive wave interference and the principle of superposition.
When two identical waves arrive at the same point exactly in phase the crests of the two waves are precisely aligned, as are the
troughs. This superposition produces pure constructive interference. Because the disturbances add, constructive interference may
produce a wave that has twice the amplitude of the individual waves, but has the same wavelength.
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Constructive Interference: Pure constructive interference of two identical waves produces one with twice the amplitude, but the
same wavelength.
If two identical waves that arrive exactly out of phase—that is, precisely aligned crest to trough—they may produce pure
destructive interference. Because the disturbances are in the opposite direction for this superposition, the resulting amplitude may
be zero for destructive interference, and the waves completely cancel.
While pure constructive and pure destructive interference do occur, they require precisely aligned identical waves. The
superposition of most waves produces a combination of constructive and destructive interference and can vary from place to place
and time to time. Here again, the disturbances add and subtract, producing a more complicated looking wave.
Superposition of Non-Identical Waves: Superposition of non-identical waves exhibits both constructive and destructive
interference.
learning objectives
Describe properties of a standing wave
Standing Wave
Sometimes waves do not seem to move, but rather they just vibrate in place. These waves are formed by the superposition of two or
more moving waves for two identical waves moving in opposite directions. The waves move through each other with their
disturbances adding as they go by. If the two waves have the same amplitude and wavelength then they alternate between
constructive and destructive interference. The resultant looks like a wave standing in place and, thus, is called a standing wave.
Standing Wave: A standing wave (black) depicted as the sum of two propagating waves traveling in opposite directions (red and
blue).
Standing waves are found on the strings of musical instruments and are due to reflections of waves from the ends of the string.
shows seven standing waves that can be created on a string that is fixed at both ends. Nodes are the points where the string does not
move; more generally, nodes are where the wave disturbance is zero in a standing wave. The fixed ends of strings must be nodes,
too, because the string cannot move there. The word antinode is used to denote the location of maximum amplitude in standing
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waves. Standing waves on strings have a frequency that is related to the propagation speed vw of the disturbance on the string. The
wavelength λ is determined by the distance between the points where the string is fixed in place.
Standing Waves: Standing waves in a string, the fundamental mode and the first six overtones.
The lowest frequency, called the fundamental frequency, is thus for the longest wavelength, twice the length of the string. The
overtones or harmonics are multiples of the fundamental frequency. shows the fundamental mode along with six overtones.
Resonance
A closer look at earthquakes provides evidence for conditions appropriate for resonance: standing waves, and constructive and
destructive interference. A building may be vibrated for several seconds with a driving frequency matching that of the natural
frequency of the vibration of the building—producing a resonance resulting in one building collapsing while neighboring buildings
do not. Often buildings of a certain height are devastated while other taller buildings remain intact. The building height matches the
condition for setting up a standing wave for that particular height. As the earthquake waves travel along the surface of Earth and
reflect off denser rocks, constructive interference occurs at certain points. Often areas closer to the epicenter are not damaged while
areas farther away are damaged.
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High School Physics - Resonance
learning objectives
Express relationship between the wave number and the wavelength, and frequency and period, of the harmonic wave
function
In this Atom we shall consider wave motion resulting from harmonic vibrations and discuss harmonic transverse wave in the
context of a string. We assume there is no loss of energy during transmission of wave along the string. This can be approximated
when the string is light and taught. In such condition, if we oscillate the free end in harmonic manner, then the vibrations in the
string are simple harmonic motion (SHM), perpendicular to the direction of wave motion. The amplitude of wave form remains
intact through its passage along the string.
We know that a traveling wave function representing motion in x-direction has the form:
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(Read our Atom on “Mathematical Representation of a Traveling Wave. “) For the case of harmonic vibration, we represent
harmonic wave motion in terms of either harmonic sine or cosine function:
y(x, t) = A sin(kx − ωt). (15.6.12)
Harmonic Waves: Harmonic waves are described by sinusoidal functions. The wavelength is equal to linear distance between
repetitions of transverse disturbance or phase.
Clearly, the displacement in y-direction is described by the bounded sine or cosine function. The important point here is to realize
that oscillatory attributes (like time period, angular and linear frequency) of wave motion is same as that of vibration of a particle in
transverse direction.
We know that time period in SHM is equal to time taken by the particle to complete one oscillation. It means that displacement of
the particle from the mean position at a given position such as x=0 has same value after time period “T” for:
ωT = 2π . Therefore, ω = 2π
T
.
Similarly, displacement of the particle from the mean position at a given time such as t=0 has same value by change the position by
“, where kλ = 2π . k is called wavenumber.
We can determine speed of the wave by noting that wave travels a linear distance ” in one period (T). Thus, speed of wave is given
by:
λ ω
v = = . (15.6.14)
T k
Refraction
Refraction is a surface phenomenon that occurs as the change in direction of a wave due to a change in its medium.
learning objectives
Formulate the law of conservation of energy and momentum as it is applied to refraction
Refraction is the change in direction of a wave due to a change in its medium. Essentially, it is a surface phenomenon—mainly in
governance to the law of conservation of energy and momentum. Due to change of medium, the phase velocity of the wave is
changed but its frequency remains constant (most commonly observed when a wave passes from one medium to another at any
angle other than 90° or 0°). Refraction of light is the most commonly observed phenomenon, but any type of wave can refract when
it interacts with a medium (e.g., when sound waves pass from one medium into another or when water waves move into water of a
different depth). Refraction is described by Snell’s law, which states that for a given pair of media and a wave with a single
frequency, the ratio of the sines of the angle of incidence θ1 and angle of refraction θ2 is equivalent to the ratio of phase velocities
(v1/v2) in the two media, or equivalently, to the opposite ratio of the indices of refraction (n2/n1):
sin θ1 v1 n2
= = . (15.6.15)
sin θ2 v2 n1
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In optics, refraction is a phenomenon that often occurs when waves travel from a medium with a given refractive index to a
medium with another at an oblique angle. For example, a light ray will refract as it enters and leaves glass, assuming there is a
change in refractive index. A ray traveling along the normal (perpendicular to the boundary) will change speed, but not direction.
Refraction still occurs in this case. Understanding of refraction led to the invention of lenses and the refracting telescope.
Refraction can be seen when looking into a bowl of water, as illustrated in. Air has a refractive index of about 1.0003, and water
has a refractive index of about 1.33. If a person looks at a straight object, such as a pencil or straw, placed partially in the water at a
slant, the object appears to bend at the water’s surface. This is due to the bending of light rays as they move from the water to the
air. Once the rays reach the eye, the eye traces them back as straight lines (lines of sight). The lines of sight (shown as dashed lines)
intersect at a higher position than where the actual rays originated (causing the pencil to appear higher and the water to appear
shallower than they actually are).
Refraction in Water: An object (in this case a pencil) partially immersed in water looks bent due to refraction: the light waves
from X change direction and so seem to originate at Y. (More accurately, for any angle of view, Y should be vertically above X,
and the pencil should appear shorter, not longer as shown. )
Diffraction
Diffraction refers to various phenomena such as the bending of waves around obstacles and the spreading out of waves past small
openings.
learning objectives
Describe the phenomenon of diffraction, according to classical physics
Diffraction refers to various phenomena which occur when a wave encounters an obstacle. In classical physics, the diffraction
phenomenon is described as the apparent bending of waves around small obstacles and the spreading out of waves past small
openings. Similar effects occur when a light wave travels through a medium with a varying refractive index, or a sound wave
travels through one with varying acoustic impedance.
Diffraction occurs with all waves, including sound waves, water waves, and electromagnetic waves such as visible light, X-rays
and radio waves. As physical objects have wave-like properties (at the atomic level), diffraction also occurs with matter and can be
studied according to the principles of quantum mechanics.
Diffraction effects are generally most pronounced for waves whose wavelengths are roughly similar to the dimensions of the
diffracting objects. If the obstructing object provides multiple, closely-spaced openings, a complex pattern of varying intensity can
result. This is due to the superposition, or interference, of different parts of a wave that travel to the observer by different paths. A
good example would be diffraction gratings.
15.6.8 https://phys.libretexts.org/@go/page/14534
Intensity Pattern: Intensity pattern formed on a screen by diffraction from a square aperture.
The effects of diffraction are often seen in everyday life. The most striking examples of diffraction are those involving light. For
example, the closely spaced tracks on a CD or DVD act as a diffraction grating to form the familiar rainbow pattern seen when
looking at a disk. This principle can be extended to engineer a grating with a structure such that it will produce any diffraction
pattern desired, like the hologram on a credit card. Diffraction in the atmosphere by small particles can cause a bright ring to be
visible around a bright light source like the sun or the moon. A shadow of a solid object, using light from a compact source, shows
small fringes near its edges. The speckle pattern which is observed when laser light falls on an optically rough surface is also a
diffraction phenomenon. All these effects are a consequence of the fact that light propagates as a wave.
∂t
2
=c
2 ∂ u
∂x
2
is given as u(x, t) = f (x + ct) + g(x − ct) , where f and g are
arbitrary functions.
learning objectives
Formulate solution of the wave equation for a traveling wave
is a solution of the wave equation for c = . In this atom, we will obtain a general mathematical form of a traveling wave.
ω
In the middle, we used the equation 1 along with the fact that partial derivatives are interchangeable.
To solve the equation 1, let’s introduce new variables: ϕ = x − ct, ψ = x + ct. From the chain rules,
15.6.9 https://phys.libretexts.org/@go/page/14534
∂ ∂ϕ ∂ ∂ψ ∂ ∂ ∂
= + = −c +c . (15.6.21)
∂t ∂t ∂ϕ ∂t ∂ψ ∂ϕ ∂ψ
∂ ∂ϕ ∂ ∂ψ ∂ ∂ ∂
= + = + . (15.6.22)
∂x ∂x ∂ϕ ∂x ∂ψ ∂ϕ ∂ψ
∂u+ ∂u−
With the change of variables, the equation 1 becomes ∂ϕ
=0 for the equation with the “+” sign and ∂ψ
=0 for the “-” sign.
Therefore, we see that
where f and g are arbitrary functions. Converting back to the original variables of x and t, we conclude that the solution of the
original wave equation is
f(x+ct) represents a left-going traveling wave, while g(x-ct) represents a right-going traveling wave. In other words, solutions of the
1D wave equation are sums of a left traveling function f and a right traveling function g. “Traveling” means that the shape of these
individual arbitrary functions with respect to x stays constant, however the functions are translated left and right with time at the
speed c. This solution was derived by Jean le Rond d’Alembert.
Boundary Condition
Any function that contains “x+ct” or “x-ct” can be a solution of the wave equation. The wave function is further determined by
taking additional information, usually given as boundary conditions and some others. For example, in the case of a string in a
guitar, we know that the wave has zero amplitude at both ends: u(x=0)=u(x=L)=0. Also, the shape of the function at an instance can
be provided to determine the function.
Wave Equation in Two Dimensions: A solution of the wave equation in two dimensions with a zero-displacement boundary
condition along the entire outer edge.
learning objectives
Describe relationship between the energy and the amplitude, and energy and intensity, of a wave
All waves carry energy. This is seen in practical applications (e.g., in medicine), as well as effects in nature. Some examples of are:
ultrasound used for deep-heat treatment of muscle strains
a laser beam to burn away malignant tissue
water waves that erode beaches
earthquakes that topple cities
The amount of energy in a wave is related to its amplitude. Large-amplitude earthquakes produce large ground displacements, as
seen in. Loud sounds have higher pressure amplitudes and come from larger-amplitude source vibrations than soft sounds. Large
ocean breakers erode the shore more than small ones. More quantitatively, a wave is a displacement that is resisted by a restoring
force. The larger the displacement x, the larger the force F=-kx needed to create it. Because work W is related to force multiplied
by distance (Fx) and energy is put into the wave by the work done to create it, the energy in a wave is related to amplitude. In fact,
a wave’s energy is directly proportional to its amplitude squared because:
15.6.10 https://phys.libretexts.org/@go/page/14534
Earthquake Destruction: The destructive effect of an earthquake is palpable evidence of the energy carried in these waves. The
Richter scale rating of earthquakes is related to both their amplitude and the energy they carry.
1
2
W =∫ F(x)dx = kx . (15.6.25)
2
The energy effects of a wave depend on time as well as amplitude. For example, the longer deep-heat ultrasound is applied, the
more energy it transfers. Therefore, power is more appropriate than energy to describe the “intensity” of a wave. Waves can also be
concentrated or spread out. Sunlight, for example, can be focused to burn wood. Earthquakes “spread out” so they do less damage
the farther they spread from their source. In both cases, changing the area the waves cover has important effects. All these pertinent
factors are included in the definition of intensity I as power (P) per unit area:
P
I = , (15.6.26)
A
Key Points
Reflection is a wave phenomenon that changes the direction of a wavefront at an interface between two different media so that
the wavefront returns into the medium from which it originated.
At the boundary, a wave must be continuous and there should be no kinks in it.
By imposing boundary conditions, we can solve wave equation and get the form of the waves. Reflection and transmission
coefficients are defined as ratio of reflected/transmitted amplitudes and the incoming amplitude.
The disturbances of waves are superimposed when they come together—a phenomenon called superposition.
As a result of superposition of waves, inteference can be observed. Interference is an effect caused by two or more waves.
Waves can interfere constructively or destructively.
The superposition of most waves produces a combination of constructive and destructive interference and can vary from place
to place and time to time.
If two waves with the same amplitude and wavelength travel in opposite directions they alternate between constructive and
destructive interference. The resultant looks like a wave standing in place and, thus, is called a standing wave.
Nodes are points of no motion in standing waves. An antinode is the location of maximum amplitude of a standing wave.
During an earthquake, buildings with a certain height may collapse more easily. This occurs when the building height matches
the condition for setting up a standing wave for that particular height.
We represent harmonic wave motion in terms of either harmonic sine or cosine function: y(x, t) = A sin(kx − ωt) .
k and ω in the harmonic wave functinon are related to wavelength and period as follows: k = , ω = .
2π
λ
2π
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Refraction is described by Snell’s law, which states that for a given pair of media and a wave with a single frequency,
sin θ1 v1 n2
= = .
sin θ2 v2 n1
∂t
= ±c
∂u
∂x
is a solution the wave equation. To solve this new equation, we introduced new
variables ϕ = x − ct, ψ = x + ct .
The solutions of the 1D wave equation are sums of a left traveling function and a right traveling function.
The wave function is further determined by taking additional information, usually given as boundary conditions and some
others.
The energy effects of a wave depend on the amplitude and duration (time) of the wave. Waves can also be concentrated or
spread out. Considering all these factors, intensity is defined as power per unit area.
In the classical wave theory, energy of a wave doesn’t depend on the frequency of the wave. However, the energy of individual
photons in a beam is determined by the frequency of the beam.
Wave’s energy is directly proportional to its amplitude squared.
Key Terms
boundary condition: A set of restraints at the boundaries, used to solve a differential equation.
superposition: The summing of two or more field contributions occupying the same space.
interference: An effect caused by the superposition of two systems of waves, such as a distortion on a broadcast signal due to
atmospheric or other effects.
constructive interference: Occurs when waves interfere with each other crest to crest and the waves are exactly in phase with
each other.
destructive interference: Occurs when waves interfere with each other crest to trough (peak to valley) and are exactly out of
phase with each other.
resonance: The increase in the amplitude of an oscillation of a system under the influence of a periodic force whose frequency
is close to that of the system’s natural frequency.
simple harmonic motion: (SHM) — Oscillating motion (as of a pendulum) in which the acceleration of the oscillator has an
equal magnitude but opposite direction to the displacement of it from the equilibrium position.
Snell’s law: A formula used to describe the relationship between the angles of incidence and refraction.
refractive index: The ratio of the speed of light in air or vacuum to that in another medium.
wave equation: An important second-order linear partial differential equation for the description of waves such as sound
waves, light waves, and water waves.
Restoring force: If the system is perturbed away from the equilibrium, the restoring force will tend to bring the system back
toward equilibrium. The restoring force is a function only of position of the mass or particle. It is always directed back toward
the equilibrium position of the system.An example is the action of a spring. An idealized spring exerts a force that is
proportional to the amount of deformation of the spring from its equilibrium length, exerted in a direction to oppose the
deformation. Pulling the spring to a greater length causes it to exert a force that brings the spring back toward its equilibrium
length. The amount of force can be determined by multiplying the spring constant of the spring by the amount of stretch.
ultrasound: Sound with a frequency greater than the upper limit of human hearing; approximately 20 kilohertz.
photoelectric effects: In photoelectric effects, electrons are emitted from matter (metals and non-metallic solids, liquids or
gases) as a consequence of their absorption of energy from electromagnetic radiation.
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15.7: Waves on Strings
learning objectives
Calculate the speed of a wave on a string
When studying waves, it is helpful to use a string to observe the physical properties of waves visually. Imagine you are holding one
end of a string, and the other end is secured and the string is pulled tight. Now, if you were to flick the string either up and down.
The wave that occurs due to this motion is called a transverse wave. A transverse wave is defined as a wave where the movement
of the particles of the medium is perpendicular to the direction of the propagation of the wave. Figure 1 shows this in a diagram. In
this case, the medium through which the waves propagate is the rope. The wave traveled from one end to the other, while the rope
moved up and down.
Figure 1: In transverse waves, the media the wave is traveling in moves perpendicular to the direction of the wave.
Wave Properties
Transverse waves have what are called peaks and troughs. The peak is the crest, or top point of the wave and the trough is the
valley or bottom point of the wave. Refer to Figure 2 for a visual representation of these terms.The amplitude is the maximum
displacement of a particle from its equilibrium position.Wavelength, usually denoted with a lambda (λ) and measured in meters, is
the distance from either one peak to the next peak, or one trough to the next trough.Period, usually denoted as T and measured in
seconds, is the time it takes for two successive peaks, or one wavelength, to pass through a fixed point.Frequency, f, is the number
of wavelengths that pass through a given point in 1 second. Frequency is measured by taking the reciprocal of a period: f = 1
Figure 2: Peaks are the top most points of the waves and troughs are the bottom, or valleys of the waves.
Speed of a Wave on a String
Velocity is found by dividing the distance traveled by the time it took to travel that distance. In waves, this is found by dividing the
wavelength by the period: v=λTv=λTWe can take the inverse proportionality to period and frequency and apply it to this situation:
λ
v = (15.7.1)
T
1
v =λ (15.7.2)
T
v = λf (15.7.3)
Reflections
When transverse waves in strings meet one end, they are reflected, and when the incident wave meets the reflected wave,
interference occurs.
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learning objectives
Explain when a standing wave occurs
Overview
Imagine you are holding one end of a string, and the other end is secured and the string is pulled tight. Now, if you were to flick the
string either up and down. The wave that occurs due to this motion is called a transverse wave. A transverse wave is defined as a
wave where the movement of the particles of the medium is perpendicular to the direction of the propagation of the wave. shows
this in a diagram. In this case, the medium through which the waves propagate is the rope. The wave traveled from one end to the
other, while the rope moved up and down.
Transverse Wave: Diagram of a transverse wave. The wave motion moves perpendicular to the medium it is traveling in.
Properties of Waves
Transverse waves have what are called peaks and troughs. The peak is the crest, or top point of the wave and the trough is the
valley or bottom point of the wave.
The amplitude is the maximum displacement of a particle from its equilibrium position.
Wavelength, usually denoted with a lambda (λ) and measured in meters, is the distance from either one peak to the next peak, or
one trough to the next trough.
Period, usually denoted as T and measured in seconds, is the time it takes for two successive peaks, or one wavelength, to pass
through a fixed point.
Frequency, f, is the number of wavelengths that pass through a given point in 1 second. Frequency is measured by taking the
reciprocal of a period: f = 1
Transverse waves can occur while being fixed at the end point or while being free at the end point.
Transverse Wave With a Fixed End Point: A transverse wave that is fixed at the end point. The reflected wave is inverted.
is an image of a transverse wave on a string that meets a free end. The wave is reflected, but unlike a transverse wave with a fixed
end, it is not inverted.
Transverse Wave With a Free End: When a transverse wave meets a free end, it is reflected.
Standing Waves
When either of the two scenarios of wave reflection occurs, the incident wave meets the reflected wave. These waves move past
each other in opposite directions, causing interference. When these two waves have the same frequency, the product of this is called
the standing waves. Standing waves appear to be standing still, hence the name. To understand how standing waves occur, we can
15.7.2 https://phys.libretexts.org/@go/page/15600
analyze them further: When the incident wave and reflected wave first meet, both waves have an amplitude is zero. As the waves
continue to move past each other, they continue to interfere with each other either constructively of destructively.
As you may remember from previous atoms, when waves are completely in phase and interfere with each other constructively, they
are amplified, and when they are completely out of phase and interfere destructively they cancel out. As the waves continue to
move past each other, and are reflected from the opposite end, they continue to interfere both ways, and a standing wave is
produced.
Every point in the medium containing a standing wave oscillates up and down and the amplitude of the oscillations depends on the
location of the point. When we observe standing waves on strings, it looks like the wave is not moving and standing still. The
principle of standing waves is the basis of resonance and how many musical instruments get their sound. The points in a standing
wave that appear to remain flat and do not move are called nodes. The points which reach the maximum oscillation height are
called antinodes.
Key Points
The type of wave that occurs in a string is called a transverse wave. In a transverse wave, the wave direction is perpendicular
the the direction that the string oscillates in.
The period of a wave is indirectly proportional to the frequency of the wave: T = .
1
The speed of a wave is proportional to the wavelength and indirectly proportional to the period of the wave: v = .λ
This equation can be simplified by using the relationship between frequency and period:v = λf .
When a transverse wave on a string is fixed at the end point, the reflected wave is inverted from the incident wave. When a
transverse wave on a string is free at the end point, the reflected wave is not inverted from the incident wave.
A standing wave occurs when an incident wave meets a reflected wave on a string.
The points in a standing wave that appear to remain flat and do not move are called nodes. The points which reach the
maximum oscillation height are called antinodes.
Every point in the medium containing a standing wave oscillates up and down and the amplitude of the oscillations depends on
the location of the point.
A standing wave has some points that remain flat due to destructive interference. These are called antinodes.
The points on a standing wave that have reached maximum oscillation do so from constructive interference, and are called
nodes.
Key Terms
transverse wave: Any wave in which the direction of disturbance is perpendicular to the direction of travel.
oscillate: To swing back and forth, especially if with a regular rhythm.
amplitude: The maximum absolute value of some quantity that varies.
standing wave: A wave form which occurs in a limited, fixed medium in such a way that the reflected wave coincides with the
produced wave. A common example is the vibration of the strings on a musical stringed instrument.
transverse wave: Any wave in which the direction of disturbance is perpendicular to the direction of travel.
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CHAPTER OVERVIEW
16: Sound
Topic hierarchy
16.1: Introduction
16.2: Sound Intensity and Level
16.3: Doppler Effect and Sonic Booms
16.4: Interactions with Sound Waves
16.5: Further Topics
This page titled 16: Sound is shared under a not declared license and was authored, remixed, and/or curated by Boundless.
1
16.1: Introduction
learning objectives
Explain how humans can characterize sound
Before delving too far into the physics of calculating sound waves (studying things like the Doppler Effect, for example), lets
address some basics of sound. Sound is a wave—a longitudinal wave of pressure that travels through compressible medias (i.e.,
solid, liquid, gaseous, or made of plasma ). There is no sound in a vacuum; by definition, a vacuum is a space free of any particles
or matter. Thus there in a vacuum, there is no media through which sound waves can travel. Following are some characteristics of
sound:
Sound travels in longitudinal waves. When drawn these are also called sinusoidal waves, a visual example of which is shown in
(we will cover this in more detail in a different section).
Sound waves have frequency; that is, the pitch of sounds goes up or down.
The amplitude of a sound determines its volume (loudness).
Tone is a measure of the quality of a sound wave.
Sound travels faster in a hot meduim, or in a solid. It also travels faster at sea level (where air pressure is higher).
Sound intensity is the energy transmitted over a certain area. Intensity is a measure of the sound’s frequency.
Ultrasound uses sound waves with high frequencies to see things normally hard to detect, like tumors. Animals, like bats and
dolphins, use ultrasound (echolocation) to navigate and locate things. Ships also use a similar technique (known as SONAR) to
locate things underwater. (This point will be discussed further in a more advanced Atom. )
Sound Perception
Every sound wave has properties that define its frequency, such as wavelength, amplitude and intensity. Calculating these
properties is outside the scope of this atom and will be addressed later. For now, it is important to know the basics of sound. As
with light waves, sound frequencies have a range. Each living creature has a different level of sound perception. For example,
consider the following examples of sound ranges (in Hz, Hertz):
Humans 20 – 20,000 Hz
Dogs 50 – 45,000 Hz
Bats 20 – 120,000 Hz
By this comparison, humans have a relatively low sound perception.
Sound Speed
Breaking the Sound Barrier: This familiar image is of a plane that is moving faster than the speed of sound.
As mentioned previously, the speed at which sound travels depends on the media through which the sound is traveling. It is much
faster in a solid than in a liquid or gas. The general formula for calculating the speed of sound is given as:
−−
K
c =√ , (16.1.1)
ρ
where K is the coefficient of stiffness of the material (also called the Bulk modulus) and p is the density of the material. We will
examine this further in another section. Generally, the expression ‘faster than the speed of sound’ refers to 344 m/s. is an image
demonstrating a plane moving faster than the speed of sound. This general measurement is taken at sea level—at a temperature of
21 degrees Celsius under normal atmospheric conditions.
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Frequency of Sound Waves
Frequency is the number of occurrences of a repeating event per unit of time. The perception of frequency is called pitch.
learning objectives
Relate frequency with the wavelength and speed of sound
Sound waves, like all other waves, have a property called frequency. Frequency is the number of occurrences of a repeating event
per unit of time. The perception of frequency is called pitch.
vs
Frequency is dependent on wavelength and the speed of sound. It is calculated with the following equation: f = λ
Frequency: A sound wave emanates from a source vibrating at a frequency f, propagates at v, and has a wavelength λ.
Alternatively, you can use the frequency and the wavelength to find the speed of sound in a specific medium. Remember that sound
travels at different speeds in different mediums; sound moves fastest through a solid. The following equation is used to find the
specific speed of sound, and is often easier to use than the standard speed of sound equation: vs=f∗λv_s=f*\lambda
A period is the duration of one cycle of a repeating event, and is the reciprocal or inverse of the frequency. The following
animation shows different frequencies and their periods, from lowest to highest.
Frequency Animation: Three flashing lights, from lowest frequency (top) to highest frequency (bottom). f is the frequency in hertz
(Hz); or the number of cycles per second. T is the period in seconds (s); or the number of seconds per cycle. T and f are reciprocals.
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Hertz
The SI unit of frequency is called a Hertz, denoted Hz. A hertz is defined as the number of cycles per second. For example, 100 Hz
signifies 100 cycles per second.
Different species can hear different frequency ranges. Humans can only hear from 20 Hz to 20,000 Hz, while dogs can hear up to
60,000 Hz. Bats can hear the highest ranges, up to 120,000 Hz. Bats use this super hearing, or ultrasound, to locate objects and
prey. By bouncing sound waves off of another object and hearing how long it takes for the sound to echo back to them, they are
able to approximate the distance between themselves and the object. This is called echolocation.
learning objectives
Calculate the frequency of the sound wave produced by the string and a column of air
Sound can be produced in many ways, by many different instruments and devices. A vibrating string or air column can both make
interesting music and it is good to understand the physics behind both.
Vibrating Strings
There are many instruments that produce sound based on strings. Guitars, cellos, pianos and many other examples. These sounds
are produced by standing waves in the strings. These waves and their frequencies are constant, and therefor the sound and pitch
produced by them is constant. This figure shows a visual of a standing wave in a string: The speed of the wave is proportional to
the root of the string tension, and inversely proportional to the root of the string density, shown by the following equation
−
−
v =√
T
μ
Pitch, and the way the sound is perceived depends on the frequency of the sound wave. Using the wave velocity, you can
find the frequency using this equation: f = v
2L
, where L is the length of the string.
Standing Waves in a String: Vibration, standing waves in a string. The fundamental and the first 6 overtones form a harmonic
series.
Air Columns
Many instruments make music by causing a vibration to a column of air in a resonator. These are usually called wind instruments.
Examples of these are flutes, recorders, and saxophones. They are either open at both ends, or closed at one end and open at the
other. The pitch of the sound is a function of the frequency. There are many factors that go into finding the frequency. Lets start
with the tubes that are closed at one end and open at the other.
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Closed Air Tubes:The maximum displacement of the air occurs at the open end of the tube, and is called the antinode. The air
movement is constrained at the closed end, and there is no displacement, and this is called the node. The distance from the node to
the antinode is 1/4 the length of one wavelength, and equal to the length of the tube, as shown in this equation:λ = 4L This can also
vw vw
be seen in this figure: The frequency is equal to the speed of sound in the air divided by the wavelength, or:f = λ
=
4L
, where
vw is the speed of sound in the air, which we learned how to find in a previous atom.
An Air column in a tube closed at one end: The fundamental and three lowest overtones for a tube closed at one end. All have
maximum air displacements at the open end and none at the closed end.
Open Air tubes:Air tubes can also be open at both ends. They are very similar to the ones we talked about above, but there is an
antinode at both ends, since they are both open, as shown in this figure: Since there is an antinode at both ends, we can see that the
vw vw
length of a wavelength is found by this equation: λ = 2L And the frequency can be found in the following equation:f = λ
=
2L
Quality of Sound
Sound quality is an assessment of accuracy or enjoyability of how a sound is perceived.
Skills to Develop
Identify factors that influence the sound quality of a reproduction of music
Sound quality is an assessment of accuracy or enjoyability of how a sound is perceived. It can be measured objectively or
subjectively. What is the difference? Objective measurement is taken when tools are used to gauge accuracy. Subjective
measurement is more of an opinion. When human listeners hear a sound and compare it to another sound they have heard, and
decide which one they enjoy more, this is a subjective measurement of sound quality.
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Reproduced Sound Quality
The sound quality of a reproduction of music depends on many factors. These include:
Recording equipment
Processing and mastering
Reproduction equipment
Listening environment
When you listen to a recording of music, you are hearing not only the music itself, but also any interference that may have been
mistakenly recorded as well. The simplest form of digitally-stored audio is PCM.
As samples are placed closer together in time, higher frequencies can be reproduced. According to the sampling theorem, any
signal with bandwidth B can be perfectly described by more than 2B samples per second. Audio must be sampled at above 40kHz:
44.1kHz for CD recordings and 48kHz for DVD recordings.
The amount of space required to store PCM depends on the number of bits per sample, the number of samples per second, and the
number of channels. For CD audio, this is 44,100 samples per second, 16 bits per sample, and 2 channels for stereo audio leading to
1,411,200 bits per second.
Speed of Sound
The speed of sound is is the distance traveled in a unit of time by a sound wave through an elastic medium, and is usually given as
344 m/s.
Skills to Develop
Calculate the speed of sound from the properties of the media
Speed of Soud
The speed of sound is is the distance traveled in a unit of time by a sound wave through an elastic medium. This medium can be a
solid, liquid, gas or even plasma. The speed of sound is dependent on the properties of the media the sound is travelling through.
The general value given for the speed of sound is the speed of a sound wave in air, at sea level, at normal atmospheric pressure; that
number is 344 m/s. However, this number is not constant. Sound travels faster in a solid than in a liquid, and faster in a liquid that
in a gas.
s
×√
T
273K
T is the temperature in Kelvin.
Mach Number
You may have heard the term Mach number in relation to speed of space craft or jets before. This is a ratio of an object’s speed in
relation to the speed of sound. The Mach number is given by the following, dimensionless equation: M = – Mach numberv –
v
Velocity of object a – Speed of sound in medium. If something is travelling at the speed of sound, that would make the equation
equal to 1, and can be denoted as Mach 1. shows a jet that is travelling at the speed of sound or faster. The vapor cone is made just
before it reaches the speed of sound and is caused by a sudden drop in air pressure.
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Faster than the Speed of Sound: This is a jet that is just about to break the sound barrier.
Key Points
Sound travels in longitudinal sinusoidal waves.
Humans can characterize sound by frequency, amplitude and tone.
The speed at which sound travels depends on the media through which the sound is traveling. It can be calculated using:
−
−
c =√
K
ρ
, where K is the coefficient of stiffness of the material (also called the Bulk modulus) and p is the density of the
material.
vs
Frequency is dependent on wavelength and the speed of sound. It is calculated with the following equation: f = λ
A period is the duration of one cycle of a repeating event, and is the reciprocal or inverse of the frequency.
The SI unit of frequency is called a Hertz, denoted Hz. A hertz is defined as the number of cycles per second. For example, 100
Hz signifies 100 cycles per second.
The frequency, and therefore the pitch of a string instrument depends on the velocity of the sound wave and the length of the
T
√
μ
Objective measurement is taken when tools are used to gauge accuracy. Subjective measurement is more of an opinion. When
human listeners hear a sound and compare it to another sound they have heard, and decide which one they enjoy more, this is a
subjective measurement of sound quality.
The quality of sound of live music is referred to as its tone. When you hear a note, you mostly hear the fundamental frequency.
But there are other harmonics present. You hear these, but they are much, much fainter than the fundamental, or main
frequency. These are called overtones.
The sound quality of a reproduction of music depends on many factors. These include: recording equipment, processing and
mastering,reproduction equipment, and even listening environment.
Sound can travel through any compressible material. These media can be solid, liquid, gas, or even plasma.
The speed of sound is dependent on the properties of the material it travels through. It will travel faster through a solid than a
liquid, and faster through a liquid than a gas.
The general number given for the speed of sound is calculated at sea level, in air, at normal atmospheric pressure. That value is
344 m/s.
Key Terms
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time t in which it occurs: f = . n
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Boundless. Provided by: Boundless Learning. Located at: www.boundless.com//physics/definition/node. License: CC BY-
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elasticity. Provided by: Wiktionary. Located at: en.wiktionary.org/wiki/elasticity. License: CC BY-SA: Attribution-
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16.2: Sound Intensity and Level
learning objectives
Calculate sound intensity from power
Overview of Intensity
Sound Intensity is the power per unit area carried by a wave. Power is the rate that energy is transferred by a wave.
is watts per meter squared or,Wm2Wm2. This is the general intensity formula, but lets look at it from a sound perspective.
Sound Intensity
2
Δp
Sound intensity can be found from the following equation: I =
2ρvw
. Δp – change in pressure, or amplitudeρ – density of the
material the sound is traveling throughvw – speed of observed sound.Now we have a way to calculate the sound intensity, so lets
talk about observed intensity. The pressure variation, amplitude, is proportional to the intensity, So it is safe to say that the larger
your sound wave oscillation, the more intense your sound will be. This figure shows this concept.
16.2.1 https://phys.libretexts.org/@go/page/14537
Sound Intensity: Graphs of the gauge pressures in two sound waves of different intensities. The more intense sound is produced
by a source that has larger-amplitude oscillations and has greater pressure maxima and minima. Because pressures are higher in the
greater-intensity sound, it can exert larger forces on the objects it encounters
Although the units for sound intensity are technically watts per meter squared, it is much more common for it to be referred to as
decibels, dB. A decibel is a ratio of the observed amplitude, or intensity level to a reference, which is 0 dB. The equation for this is:
\(\mathrm{β=10 \log _{10} \frac{I}{I_0}β}\) – decibel levelI – Observed intensityI0– Reference intensity.For more on decibels,
please refer to the Decibel Atom.
For a reference point on intensity levels, below are a list of a few different intensities:
0 dB, I = 1×10-12 –> Threshold of human hearing
10 dB, I = 1×10-11 –> Rustle of leaves
60 dB, I = 1×10-6 –> Normal conversation
100 dB, I = 1×10-2 –> Loud siren
160 dB, I = 1×104–> You just burst your eardrums
Skills to Develop
Explain how frequency is perceived by humans
The study of the human perception of sound is called psychoacoustics. Many factors go into hearing, including wave properties,
sensory and brain processes. First, the wave has to be made, and it has a specific wavelength and frequency. Then the sound wave
reaches the human ear, and is processed through many areas. Finally, the sound wave makes it through the ear and to the human
brain, where even more action happens. You might think that when something makes a noise that you hear it instantaneously but, in
reality, it goes through many steps first.
Wave Properties
We are not going to go into too much detail about the wave’s physical properties, since it is out of the scope of this atom, but
remember:
Frequency is perceived by humans as pitch;
The sound intensity is the amplitude;
Humans can only hear a specific range of sound, usually from 20 Hz to 20,000 Hz;
The factors that go into a sound are its intensity, frequency and overtones (which are like interference, or background noises).
16.2.2 https://phys.libretexts.org/@go/page/14537
The Human Ear
The human ear is made up of three main sections, as shown in:
Decibels
The decibel is a logarithmic unit used to quantify sound levels, by comparing a physical quantity to a reference level.
Skills to Develop
Identify how decibel is quantified in acoustics
The decibel, dB, is commonly used to quantify sound levels, although it is not a unit of sound, but a unit of pressure. The decibel is
a logarithmic unit that indicates the ratio if a physical quantity to a reference level. It is one tenth of a Bel, which was named after
the inventor of the telephone, Alexander Graham Bell. The word decibell comes from the prefix, deci, that is 1/10 of the word it
precedes For more information on how to convert units, refer to the unit conversion atom. Although the decibel can be used to talk
about a number of different subjects, in this atom we are going to cover its use in acoustics and sound level.
In acoustics, the decibel is quantified relative to a reference which has been set at a sound pressure level of 20 micropascals, and is
called a 0 dB. This reference level is a typical threshold of human hearing perception. The following equation is used to calculate
the sound pressure level, or amplitude: amplitude[dB] = 20 log 10
so is the reference pressure which is 20 micropascals or 0 dB,
s
so
and s is the observed sound pressure.The human ear has a standard sound threshold of 120 dB, which expressed logarithmically is
around 1012. This is a standard threshold, but it also depends on frequency. Loudness is a measure of sound intensity taking
frequency into account, and is called a A-weighted decibel, dB(A), or a phon. This figure shows The Fletcher Munson Chart, which
demonstrates the different sound frequencies and decibels that the human ear perceives as the same.
16.2.3 https://phys.libretexts.org/@go/page/14537
The Fletcher Munson Chart: The Fletcher-Munson equal-loudness contours. Phons are labelled in blue
Key Points
2
Δp
Sound intensity can be found from the following equation: I = 2ρvw
. Δ p – change in pressure, or amplitude ρ – density of the
material the sound is traveling through vw– speed of observed sound.
The larger your sound wave oscillation, the more intense your sound will be.
Although the units for sound intensity are technically watts per meter squared, it is much more common for it to be referred to
as decibels, dB.
Frequency is perceived by humans as pitch. The sound intensity is what humans can hear, and is generally only a specific range
of sound, usually from 20 Hz to 20,000 Hz. The factors that go into a sound are its intensity, frequency and overtones (which
are like interference or background noises).
Your ear is made up of three major sections: the inner, middle and outer ear.
Your cochlea, which is in your inner ear, not only transmits sound waves to your brain, but also contains a liquid that helps
humans maintain their balance.
In acoustics, the decibel is quantified relative to a reference which has been set at a sound pressure level of 20 micropascals,
and is called a 0 dB.
The following equation is used to calculate the sound pressure level, or amplitude: amplitude[dB] = 20 log 10
. s is the
s
so
o
reference pressure which is 20 micropascals or 0 dB, and s is the observed sound pressure.
The human ear has a standard sound threshold of 120 dB, which expressed logarithmically is around 1012. This is a standard
threshold, but it also depends on frequency. Loudness is a measure of sound intensity taking frequency into account, and is
called a A-weighted decibel, dB(A), or a phon.
Key Terms
decibel: A common measure of sound intensity that is one tenth of a bel on the logarithmic intensity scale. It is defined as
dB = 10 log (P / P ) , where P and P are the relative powers of the sound.
10 1 2 1 2
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Sound Intensity and Decibels. Located at: http://www.youtube.com/watch?v=R5rkg8mTRBI. License: Public Domain: No
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Physics Study Guide/Sound. Provided by: Wikibooks. Located at: en.wikibooks.org/wiki/Physics_Study_Guide/Sound.
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phon. Provided by: Wiktionary. Located at: en.wiktionary.org/wiki/phon. License: CC BY-SA: Attribution-ShareAlike
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16.3: Doppler Effect and Sonic Booms
learning objectives
Identify parameters required to calculate the frequency perceived by the observer moving towards the sound source
In this atom, we are going to cover the Doppler effect, but specifically when the observer is the one in motion.
Sound and the Doppler Effect: This video introduces sound waves. The first video describes the basics of sound while the second
video looks at the Doppler Effect.
16.3.1 https://phys.libretexts.org/@go/page/14538
The Doppler Effect: The same sound source is radiating sound waves at a constant frequency in the same medium. However, now
the sound source is moving to the right with a speed υs = 0.7 c (Mach 0.7). The wave-fronts are produced with the same frequency
as before. However, since the source is moving, the centre of each new wavefront is now slightly displaced to the right. As a result,
the wave-fronts begin to bunch up on the right side (in front of) and spread further apart on the left side (behind) of the source.
When the observer moves toward an sound source, each successive wave is encountered sooner than the previous wave. Thus, it
will take just a little less time for the observer to hear the next one. Since the time between waves is reduced, the frequency is
increased. Similarly if the observer is moving away from the sound source, the frequency, and therefor pitch, is decreased. While
the frequency will change whether the observer or sound source is moving, it is easier to show with the sound source as the one
moving. This figure demonstrated the sound source moving:
Unless the observer is moving directly towards the sound source, this angle needs to be taken into account when calculating the
newly perceived frequency. Before we can start this calculation, we must know:
The original sound wave frequency, f0
The velocity of the observer, vr
The speed of sound in the air, or medium, c
The angle of the line of sight from the observer to the sound source, θθ
Although the sound waves are being emitted from the sound source at a uniform frequency, the observer is perceiving them
c±vr ,new
differently. The equation for the perceived wave frequency is as follows: f = c
× f And the equation for vr,
0 is:
vr,new = v cos θ If the observer is moving towards the sound source, you are going to use a plus sign in front of the observers
r
velocity. If the observer is moving away from sound source, you are going to use a negative sign in front of the observers velocity.
Moving Source
The Doppler effect is the apparent change in frequency of a wave when the observer and the source of the wave move relative to
each other.
learning objectives
Identify parameters required to calculate the perceived frequency of the sound source moving towards an observer
The Doppler Effect: The same sound source is radiating sound waves at a constant frequency in the same medium. However, now
the sound source is moving to the right with a speed υs = 0.7 c ( Mach 0.7). The wave-fronts are produced with the same frequency
16.3.2 https://phys.libretexts.org/@go/page/14538
as before. However, since the source is moving, the centre of each new wavefront is now slightly displaced to the right. As a result,
the wave-fronts begin to bunch up on the right side (in front of) and spread further apart on the left side (behind) of the source.
When the sound source moves toward an observer, each successive wave is emitted closer to the observer than the previous wave
and takes just a little less time to reach the observer than the previous one. Since the time between waves is reduced, the frequency
is increased. Similarly, if the sound source is moving away from the observer, the frequency (and therefore pitch) is decreased.
While the frequency will change whether the observer or sound source is moving, the effect is more easily demonstrated by the
sound source. This Doppler Effect is illustrated in.
Unless the observer is moving directly towards the sound source, this angle must be considered when calculating the newly
perceived frequency. Before attempting this calculation, we must know:
The original sound wave frequency, f0
The velocity of the observer, vr
The speed of sound in the air, or medium, c
The angle of the line of sight from the observer to the sound source, θ
Although the sound waves are being emitted from the sound source at a uniform frequency, the observer is perceiving them
differently. The equation for the perceived wave frequency is as follows: f = ∗f And the equation for vs, is:
c
c±vs ,new
0
If the sound source is moving towards the observer, a plus sign is used in front of the sound source’s velocity. If the sound source is
moving away from the observer, then a negative sign is used in front of the sound source’s velocity.
General Case
The Doppler effect is the apparent change in frequency of a wave when the observer and the source of the wave move relative to
each other.
learning objectives
Compare change in the frequency due to the doppler effect when the object in motion moves towards or away from the
other object
Have you ever wondered why as a firetruck speeds by, the sound of the siren seems to change? Or where the very distinguishable
sound at a race track comes from as the cars whisk by you? Well believe it or not, this is science! The name of this is called the
Doppler effect, named for the scientist who discovered it. His name was Christian Doppler, and he discovered it around 1840 in
Prague. He preformed the experiment in two segments. In the first experiment, he had a band play music while aboard a moving
train, and would have observers listen to the music in a stationary location while the train went by. In the second experiment, he
had a group of observers sit aboard a moving train, while the band played in a stationary location. In both cases, observers noticed
the phenomenon. The change in sound perception can be explained through relativity.
When the sound source moves toward an observer, each successive wave is emitted closer to the observer than the previous wave.
Thus, it will take just a little less time to reach the observer than the previous one. Since the time between waves is reduced, the
frequency is increased. Similarly if the sound source is moving away from the observer, the frequency, and therefor pitch, is
decreased. If the observer is moving away from the sound source, the frequency will be lowered, and if the observer moves closer
to the sound source, the frequency is increased. While the frequency will change whether the observer or sound source is moving, it
is easier to show with the sound source as the one moving. This figure demonstrated the sound source moving:
16.3.3 https://phys.libretexts.org/@go/page/14538
The Doppler Effect: The same sound source is radiating sound waves at a constant frequency in the same medium. However, now
the sound source is moving to the right with a speed υs = 0.7 c (Mach 0.7). The wave-fronts are produced with the same frequency
as before. However, since the source is moving, the centre of each new wavefront is now slightly displaced to the right. As a result,
the wave-fronts begin to bunch up on the right side (in front of) and spread further apart on the left side (behind) of the source.
In classical physics, where the speeds of source and the receiver relative to the medium are lower than the velocity of waves in the
(c+vr)
medium, the relationship between observed frequency (f) and emitted frequency (fo) is given by: f = (c+vs)
× f0 where c – velocity
of the sound waves in the medium, vr– velocity of the observer or reciever, vs – velocity of the sound source, and f0 – original
frequency of the sound waves.
The above formula assumes that the source is either directly approaching or receding from the observer. If the source approaches
the observer at an angle (but still with a constant velocity), the observed frequency that is first heard is higher than the object’s
emitted frequency. Thereafter, there is a monotonic decrease in the observed frequency as it gets closer to the observer, through
equality when it is closest to the observer, and a continued monotonic decrease as it recedes from the observer. When the observer
is very close to the path of the object, the transition from high to low frequency is very abrupt. When the observer is far from the
path of the object, the transition from high to low frequency is gradual.
If the speeds and are small compared to the speed of the wave, the relationship between observed frequency and emitted frequency
is approximately
Observed frequency: f = (1 + Δv
c
)f0
Change in frequency: Δf = Δv
c
f0
where Δv is the velocity of the receiver relative to the source: it is positive when the source and the receiver are moving towards
each other.
Sonic Booms
A sonic boom is the sound associated with the shock waves created by an object traveling through the air faster than the speed of
sound.
learning objectives
Identify conditions that lead to a sonic boom and discuss its properties
Sonic Booms
A sonic boom is a sound associated with the shock waves created by an object traveling through the air faster than the speed of
sound. It can be viewed as a Doppler effect on steroids; sonic booms generate an enormous amount of energy and sound like
explosions. The first man-made object to ever create this phenomenon was a bullwhip. The ‘crack’ of the whip is a result of this
sonic boom. This version of a Doppler effect is demonstrated by.
16.3.4 https://phys.libretexts.org/@go/page/14538
Sonic Boom Gif: The sound source has now broken through the sound speed barrier, and is traveling at 1.4 times the speed of
sound, (Mach 1.4). Since the source is moving faster (with a speed ) than the sound waves it creates, it actually leads the advancing
wavefront. The sound source will pass by a stationary observer (with a speed ) before the observer actually hears the sound it
creates.
When the sound source passes through the air, it creates a series of pressure waves. These waves are travelling at the speed of
sound, and as the speed of the sound source increases, the waves, not being able to get out of each other’s way, are forced together.
They eventually merge into a single shock wave traveling at the speed of sound. This is a critical speed, known as Mach. The shock
waves radiate out from the sound source, and create a “Mach cone’. The half-angle, αα, can be found using the equation
Sonic Boom: A sonic boom produced by an aircraft moving at M=2.92, calculated from the cone angle of 20 degrees. An observer
hears the boom when the shock wave, on the edges of the cone, crosses his or her location
vs
sin α = . (16.3.2)
vr
vs
From previous atoms, we know that vr
is the sound source’s Mach number.
At the front of the sound source, there is a sudden rise in pressure, while at the end of the source there is a decreasing pressure. This
‘overpressure profile’ is known as an N-wave. There is a big boom when there is a sudden change in pressure, and since the
pressure changes twice, this is a double boom.
Key Points
When the object in motion moves towards the other, the frequency is increased because the time between successive sound
waves is shortened. Therefore the pitch is higher.
When the object in motion moves away from the other, the frequency is decreased because the time between successive sound
waves is lengthened. Therefor the pitch is lowered.
Unless the objects are in each other’s direct path, you need to account for the angle they are at relative to each other. The
following equation needs to be substituted for the ‘movers’ velocity. The angle used needs to be the angle from the line of sight
of the observer to the sound source. vradial=v∗cosθvradial=v∗cosθ.
When the object in motion moves towards the other, the frequency is increased because the time between successive sound
waves is shortened (therefore the pitch is higher).
16.3.5 https://phys.libretexts.org/@go/page/14538
When the object in motion moves away from the other, the frequency is decreased because the time between successive sound
waves is lengthened (therefore the pitch is lowered).
Unless the objects are in each other’s direct path, you need to account for their angle relative to each other. The following
equation must be substituted for the ‘movers’ velocity. The angle used must be the angle from the line of sight of the observer to
the sound source.
When the object in motion moves towards the other, the frequency is increased because the time between successive sound
waves is shortened. Therefore the pitch is higher.
When the object in motion moves away from the other, the frequency is decreased because the time between successive sound
waves is lengthened. Therefor the pitch is lowered.
If the speeds and are small compared to the speed of the wave, the relationship between observed frequency and emitted
frequency is approximately: f = (1 + Δv
c
)f
0 where Δv is the velocity of the receiver relative to the source: it is positive when
the source and the receiver are moving towards each other.
A sonic boom happens when a sound source passes an observer at a speed either very close to, or faster than the speed of sound.
Sonic booms generate an enormous amount of energy, and sound like explosions.
The half angle of the cone of sound waves produced by the sonic boom can be found by taking the inverse sine of the object’s
Mach number.
Key Terms
Doppler's effect: Apparent change in frequency of a wave when the observer and the source of the wave move relative to each
other.
classical physics: All aspects of physics developed before the rise of quantum mechanics.
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time t in which it occurs: f = n/t .
Mach number: The ratio of the velocity of a body to that of sound in the surrounding medium.
sonic boom: The audible effect of a shock wave in the air, especially one caused by an aircraft flying faster than the speed of
sound
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16.4: Interactions with Sound Waves
learning objectives
Identify conditions required for the superposition of two waves
Waves are most commonly described by variations in some parameter through space and time—height in a water wave, pressure in
a sound wave, or the electromagnetic field in a light wave. The value of this parameter is called the amplitude of the wave; the
wave itself is a function specifying the amplitude at each point.
When two or more waves arrive at the same point, they superimpose themselves onto one another. More specifically, the
disturbances of waves are superimposed when they come together (a phenomenon called superposition). Each disturbance
corresponds to a force, or amplitude (and the forces add). If the disturbances are along the same line, the resulting wave is a simple
addition of the disturbances of the individual waves. That is, their amplitudes add.
Superposition: Superposition is when two waves add together. In this figure, the two waves add together and cancel out leaving no
wave. This is destructive interference.
Superposition of waves leads to what is known as interference, which manifests in two types: constructive and destructive.
Constructive interference occurs when two waves add together in superposition, creating a wave with cumulatively higher
amplitude, as shown in. In destructive interference, the two waves add together but cancel out (like adding a positive and negative
number). Destructive interference is shown in.
While pure constructive and pure destructive interference do occur, they require precisely aligned identical waves. The
superposition of most waves produces a combination of constructive and destructive interference, and can vary from place to place
and time to time. Sound from a stereo, for example, can be loud in one spot but quiet in another. Varying loudness means the sound
waves add partially constructively and partially destructively at different locations. A stereo has at least two speakers creating
sound waves, and waves can reflect from walls. All these waves superimpose. An example of sounds that vary over time from
constructive to destructive is found in the combined whine of airplane jets heard by a stationary passenger. The combined sound
can fluctuate up and down in volume as the sound from the two engines varies in time from constructive to destructive.
These examples are of waves that are similar. illustrates that when non-identical waves superimpose, the outcome is a mixture of
constructive and destructive interference.
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Superposition of Non-Identical Waves: Superposition of non-identical waves exhibits both constructive and destructive
interference.
Interference
Interference occurs when multiple waves interact with each other, and is a change in amplitude caused by several waves meeting.
learning objectives
Contrast constructive and destructive interference
Unlike solid objects, two waves can share a point in space. In physics, interference is a phenomenon in which two waves (passing
through the same point) superimpose to form a resultant wave of greater or lower amplitude. Interference usually refers to the
interaction of waves that are correlated or coherent with each other (i.e, “interfere” with each other), either because they come from
the same source or because they have the same or nearly the same frequency.
The effects of interference can be observed with all types of waves, for example, light, radio, acoustic and surface water waves.
The idea that interference is caused by superposition means that when two waves meet their two amplitudes (their maximum
absolute value) combine together.
16.4.2 https://phys.libretexts.org/@go/page/14539
Constructive Interference: Pure constructive interference of two identical waves produces one with twice the amplitude, but the
same wavelength.
Destructive interference is when two waves add together and the result is a smaller displacement than would have been the case.
An example of destructive interference can be seen in. When the waves have opposite amplitudes at the point they meet they can
destructively interfere, resulting in no amplitude at that point. For example, this is how noise cancelling headphones work. By
playing a sound with the opposite amplitude as the incoming sound, the two sound waves destructively interfere and this cancel
each other out.
Beats
The superposition of two waves of similar but not identical frequencies produces a pulsing known as a beat.
learning objectives
Identify superposition conditions that lead to beat
Striking two adjacent keys on a piano produces a warbling combination (usually considered unpleasant to the ear). The culprit is
the superposition of two waves of similar but not identical frequencies. When two waves of similar frequency arrive at the same
point and superimpose, they alternately constructively and destructively interfere. This alternating is known as a beat because it
produces an unpleasant pulsing sound.
Another example is often noticeable in a taxiing jet aircraft (particularly the two-engine variety). The loudness of the combined
sound of the engines increases and decreases. This varying loudness occurs because the sound waves have similar but not identical
frequencies. The discordant warbling of the piano and the fluctuating loudness of the jet engine noise are both due to alternately
constructive and destructive interference as the two waves go in and out of phase. illustrates this phenomenon graphically.
Beat Frequency: Beats are produced by the superposition of two waves of slightly different frequencies but identical
amplitudes.The waves alternate in time between constructive interference and destructive interference, giving the resulting wave a
time-varying amplitude.
The wave resulting from the superposition of two similar-frequency waves has a frequency that is the average of the two. This
wave fluctuates in amplitude, or beats, with a frequency called the beat frequency. We can determine the beat frequency
mathematically by adding two waves together.
One can also measure the beat frequency directly. When you hear a beat coming from two discordant sounds (say, two notes on a
piano) you can count the number of beats per second. The number of beats per second, or the beat frequency, shows the difference
in frequency between the two notes. Musicians often use this phenomena to ensure that two notes are in tune (if they are in tune
then there are no beats).
The Ear
The ear is the sensory organ that picks up sound waves from the air and turns them into nerve impulses that can be sent to the brain.
learning objectives
Describe how sound waves are collected and transformed into nerve impulses
Sound waves are vibrations in the air. The ear is the sensory organ that picks up sound waves from the surrounding air and turns
them into nerve impulses, which are then sent to the brain. The sound waves carry a lot of information — language, music, and
noise — all mixed together. The task of the ear is to turn the signals in these waves of bouncing air molecules into electrical nerve
signals while keeping as much of the information in the signal as possible. (It’s the brain’s job to then sort the signals and make
16.4.3 https://phys.libretexts.org/@go/page/14539
sense of them. ) It’s not easy to turn one kind of signal into another without losing information, but the ear is well designed for the
task.
Anatomy of the Human Ear: Anatomy of the human ear; the length of the auditory canal is exaggerated for viewing purposes
Air surrounds the head and fills the ear canal and middle ear. Therefore, when the outer part of the ear collects sound and the
middle ear amplifies this sound pressure, these processes occur in the medium of air. However, the hollow channels of the inner ear
(which is embedded in the temporal bone, the densest bone of the body) are filled with liquid. So as the sound travels into the inner
ear, it passes from the medium of air into a liquid medium. These inner-ear channels contain a sensory epithelium that is studded
with hair cells. The microscopic “hairs” are structural protein filaments that project out into the fluid. The hair cells release a
chemical neurotransmitter when stimulated. Sound waves moving through fluid push the filaments; if the filaments bend over
enough, the hair cells fire chemical signals. In this way sound waves are transformed into nerve impulses. The nerve impulses
travel from the left and right ears through the eighth cranial nerve to both sides of the brain stem and up to the part of the cerebral
cortex dedicated to sound (auditory cortex, located in the temporal lobe).
learning objectives
Discuss application of sound waves in medicine and navigation
The controlled use of sound waves has many applications in science. Controlled reflection of such waves allows images to be
received.
Ultrasound
Ultrasound is sound with a frequency higher than 20 kHz. This is above the human range of hearing. The most common use of
ultrasound, creating images, has industrial and medical applications. The use of ultrasound to create images is based on the
reflection and transmission of a wave at a boundary. When an ultrasound wave travels inside an object that is made up of different
materials (such as the human body), each time it encounters a boundary (e.g., between bone and muscle, or muscle and fat), part of
the wave is reflected and part of it is transmitted. The reflected rays are detected and used to construct an image of the object.
Sonar
illustrates how a ship on the ocean utilizes the reflecting properties of sound waves to determine the depth of the ocean. A sound
wave is transmitted and bounces off the seabed. Because the speed of sound is known and the time lapse between sending and
receiving the sound can be measured, the distance from the ship to the bottom of the ocean can be determined. This technique is
called sonar (originally an acronym for SOund Navigation And Ranging).
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Sonar: Ships on the ocean make use of the reflecting properties of sound waves to determine the depth of the ocean. A sound wave
is transmitted and bounces off the seabed. Because the speed of sound is known and the time lapse between sending and receiving
the sound can be measured, the distance from the ship to the bottom of the ocean can be determined.
Just as ships on the ocean, certain animals, like dolphins and bats, make use of sounds waves (sonar) to navigate or find their way.
Ultrasound waves are sent out then reflected off the objects around the animal. Bats or dolphins then use the reflected sounds to
form a “picture” of their surroundings (this is known as echolocation).
Key Points
When two waves occupy the same point, superposition occurs. Superposition results in adding the two waves together.
Constructive interference is when two waves superimpose and the resulting wave has a higher amplitude than the previous
waves.
Destructive interference is when two waves superimpose and cancel each other out, leading to a lower amplitude.
Most wave superpositions involve a mixture of constructive and destructive interference since the waves are not perfectly
identical.
Interference is a phenomenon of wave interactions. When two waves meet at a point, they interfere with each other.
There are two types of interference, constructive and destructive.
In constructive interference, the amplitudes of the two waves add together resulting in a higher wave at the point they meet.
In destructive interference, the two waves cancel out resulting in a lower amplitude at the point they meet.
When two waves of similar frequencies interfere, the result is a beat frequency.
A beat frequency is a pulsing sound that goes up and down in loudness.
As the two waves go in and out of phase, the varying constructive and destructive interference makes the wave grow and shrink
in amplitude. For sound waves this produces a beating sound.
The task of the ear is to turn the signals in the waves of bouncing air molecules into electrical nerve signals while keeping as
much of the information in the signals as possible.
Sound is collected in the outer part of the ear; sound pressure is amplified through the middle part of the ear and is passed from
the medium of air into a liquid medium.
That sound pressure is amplified through the middle portion of the ear and passed from the medium of air into a liquid medium.
Sound waves moving through the fluid in the inner ear stimulate hair cells, making them release chemical neurotransmitters. In
this way sound waves are transformed into nerve impulses.
When waves encounter a boundary between two materials, part of the wave is reflected and part is transmitted.
By using high frequency sound waves, doctors can create images of parts of the body normally not visible.
By transmitting sound waves and measuring the time between the transmission and receiving the reflection, ships can use sound
waves to navigate. This is called sonar.
Key Terms
superimpose: To place an object over another object.
displacement: A vector quantity that denotes distance with a directional component.
amplitude: The maximum absolute value of some quantity that varies.
coherent: Of waves having the same direction, wavelength and phase, as light in a laser.
16.4.5 https://phys.libretexts.org/@go/page/14539
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time tt in which it occurs: \
(\mathrm{f=n/t\).
interfere: (of waves) To be correlated with each other when overlapped or superposed.
superposition: The summing of two or more field contributions occupying the same space.
epithelium: a membranous tissue composed of one or more layers of cells that forms the covering of most internal and external
surfaces of the body and its organs (internally, the lining of vessels and other small cavities; externally, the skin)
nerve impulse: the signal transmitted along a nerve fiber, either in response to a stimulus (such as touch, pain, or heat), or as an
instruction (such as causing a muscle to contract)
neurotransmitter: any substance, such as acetylcholine or dopamine, responsible for sending nerve signals across a synapse
between two neurons
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16.5: Further Topics
learning objectives
Identify sources of spherical and plane wave patterns
Spherical Waves
Christiaan Huygens was a Dutch scientist who developed a useful technique for determining how and where waves propagate. In
1678, he proposed that every point that a luminous disturbance touches becomes itself a source of a spherical wave; the sum of
these secondary waves determines the form of the wave at any subsequent time. The Huygen-Fresnel Principle shows that as the
waves interact with each other, they interfere either constructively or destructively. Constructive interference occurs when waves
are completely in phase with each other and amplifies the waves. Destructive interference occurs when waves are exactly out of
phase with either other, and if waves are perfectly out of phase with each other, the wave will be canceled out completely. Since the
waves all come from one point source, the waves happen in a spherical pattern. All the waves come from a single point source and
are spherical.
Spherical Wave: When waves are produced from a point source, they are spherical waves.
Huygen-Fresnel Principle: The Huygen-Fresnel Principle uses the law of refraction. Each point on the wave produces waves
which interfere with each other either constructively or destructively.
Plane Waves
A plane wave is a constant-frequency wave whose wavefronts (surfaces of constant phase) are infinite parallel planes of constant
peak-to-peak amplitude normal to the phase velocity vector. It is not possible in practice to have a true plane wave; only a plane
wave of infinite extent will propagate as a plane wave. However, many waves are approximately plane waves in a localized region
of space. For example, a localized source such as an antenna produces a field that is approximately a plane wave far from the
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antenna in its far-field region. Similarly, if the length scales are much longer than the wave’s wavelength, as is often the case for
light in the field of optics, one can treat the waves as light rays which correspond locally to plane waves.
Plane Wave: Plane waves are an infinite number of wavefronts normal to the direction of the propogation.
learning objectives
Identify when a standing wave occurs
A standing wave is a wave that appears stationary, meaning it remains in a constant position. In a string, a standing wave is a type
of transverse wave—where the movement of the particles of the medium is perpendicular to the direction of the propagation of the
wave. A standing wave can occur when two identical waves moving in different directions along the string interfere.
There are two scenarios of waves in strings: the string is fixed at both ends, or the string is fixed at one end and free at the other. A
transverse wave will move along the string until it reaches the other end. It is then reflected from that end and starts to move back
towards the original direction; at this point interference occurs. shows a transverse wave that is reflected from a fixed end. When a
transverse wave meets a fixed end, the wave is reflected, but inverted. This swaps the peaks with the troughs and the troughs with
the peaks. diagrams a transverse wave on a string that meets a free end. The wave is reflected, but unlike a transverse wave with a
fixed end, it is not inverted.
Free End Reflection: The wave is reflected, but unlike a transverse wave with a fixed end, it is not inverted.
Fixed End Reflection: When a transverse wave meets a fixed end, the wave is reflected, but inverted.
Standing Waves
When either of the two scenarios of wave reflection occurs, the incident wave meets the reflected wave. These waves move past
each other in opposite directions, causing interference. When these two waves have the same frequency, the product of this is called
the standing waves. Standing waves appear to be standing still, hence the name. illustrates a very slow moving standing wave. (One
application of the principle of standing waves is in music with the concept of resonance—and how many musical instruments, like
guitars and pianos, get their sound. ) Let us now examine how standing waves occur.
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Standing Wave on a String: This is what a standing wave would look like if you were to slow it down. The wave is caused by an
incident wave on a string being reflected and then traveling back in the direction it came from. The two waves then meet and
interfere with each other causing this phenomenon.
learning objectives
Identify the type of a standing wave in an air column
Standing Waves
A standing wave is a wave that appears to be stationary, meaning it remains in a constant position. In an air column, a standing
wave can form as either a longitudinal or transverse wave. A transverse wave, you may remember, is a wave where the movement
of the particles of the medium is perpendicular to the direction of the propagation of the wave. A longitudinal wave, on the other
hand, is parallel to the direction of propagation. A standing wave can occur when two identical waves moving in different
directions interfere.
Air Columns
When a standing wave is formed in a tube, the standing wave has a maximum air displacement at the open end called an antinode.
Here, the motion is unconstrained. At the closed end, there is no displacement; this is called a node, and the air is halted. The
distance from a node to antinode is 1/4 of a wavelength, and is equal to the length of the tube.
shows that this resonance can be produced by a vibration introduced at or near the closed end of the tube as well. This is considered
to be a natural vibration of the air column independently of how it is induced. Given that maximum air displacements are possible
at the open end and none at the closed end, there are other, shorter wavelengths that can resonate in the tube. For example, a
standing wave can have three-fourths of its wavelength in the tube, or L=(3/4)λ. This can be rendered as follows:
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Standing Wave in Air Column: The same standing wave is created in the tube by a vibration introduced near its closed end.
4L
′
λ = (16.5.1)
3
Continuing this process reveals a whole series of shorter-wavelength and higher-frequency sounds that resonate in the tube. We use
specific terms for the resonances in any system. The lowest resonant frequency is called the fundamental, while all higher resonant
frequencies are called overtones. All resonant frequencies are integral multiples of the fundamental, and they are collectively called
harmonics. The fundamental is the first harmonic, the first overtone is the second harmonic, and so on. shows hows the
fundamental and the first three overtones (the first four harmonics) in a tube closed at one end.
The First Four Harmonics: The fundamental and three lowest overtones for a tube closed at one end. All have maximum air
displacements at the open end and none at the closed end.
Now let us look for a pattern in the resonant frequencies for a simple tube that is closed at one end. The fundamental has =4L, and
frequency is related to wavelength and the speed of sound as given by the following:
vw = f λ (16.5.2)
Here, f is frequency, vw is speed of sound in air, λ is wavelength, and L is the length of the air column. The first overtone has =
4L/3. From this, we can deduce the following:
′
vw
f =3 = 3f (16.5.4)
4L
Because f′ = 3f, we call the first overtone the third harmonic. Continuing this process, we see a pattern that can be generalized in a
single expression. The resonant frequencies of a tube closed at one end are:
vw
fn = n , n = 1, 3, 5 (16.5.5)
4L
Here, f1 is the fundamental, f3 is the first overtone, and so on. It is interesting that the resonant frequencies depend on the speed of
sound and, hence, on temperature. This dependence poses a noticeable problem for organs in old unheated cathedrals, and it is also
the reason why musicians commonly bring their wind instruments to room temperature before playing them.
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Example 16.5.1:
The fundamental and overtones can be present simultaneously in a variety of combinations. For example, middle C on a
trumpet has a sound distinctively different from middle C on a clarinet, both instruments being modified versions of a tube
closed at one end. The fundamental frequency is the same (and usually the most intense), but the overtones and their mix of
intensities are different and subject to shading by the musician. This mix is what gives various musical instruments (and human
voices) their distinctive characteristics, whether they have air columns, strings, sounding boxes, or drumheads. In fact, much of
our speech is determined by shaping the cavity formed by the throat and mouth and positioning the tongue to adjust the
fundamental and combination of overtones. Simple resonant cavities can be made to resonate with the sound of the vowels, for
example.
learning objectives
Explain the relationship between the resonance curve and damping
Paddle Balls and Frequencies: The paddle ball on its rubber band moves in response to the finger supporting it. If the finger
moves with the natural frequency f0 of the ball on the rubber band, then a resonance is achieved, and the amplitude of the ball’s
oscillations increases dramatically. At higher and lower driving frequencies, energy is transferred to the ball less efficiently, and it
responds with lower-amplitude oscillations.
As the frequency at which the finger is moved up and down increases, the ball will respond by oscillating with increasing
amplitude. After driving the ball at its natural frequency, the ball’s oscillations increase in amplitude with each oscillation for as
long as it is driven. The phenomenon of driving a system with a frequency equal to its natural frequency is called resonance. As the
driving frequency gets progressively higher than the resonant or natural frequency, the amplitude of the oscillations becomes
smaller until the oscillations nearly disappear and the finger simply moves up and down with little effect on the ball.
Effects of Damping
In real life, most oscillators have damping present in the system. It is interesting that the widths of the resonance curves shown in
depend on damping: the less the damping, the narrower the resonance. For a driven oscillator to resonate at a very specific
frequency, there needs to be as little damping as possible, as is the case for piano strings and many other musical instruments.
Conversely, for small-amplitude oscillations, such as in a car’s suspension system, there needs to be heavy damping. Heavy
damping reduces the amplitude, but the tradeoff is that the system responds at more frequencies.
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Damping: The amplitude of a harmonic oscillator is a function of the frequency of the driving force. The curves represent the same
oscillator with the same natural frequency but with different amounts of damping. Resonance occurs when the driving frequency
equals the natural frequency, and the greatest response is for the least amount of damping. The narrowest response is also for the
least amount of damping.
Example 16.5.2:
These features of driven harmonic oscillators apply to a huge variety of systems. When tuning a radio, for example, people are
adjusting the resonant frequency of the radio circuit so that it only oscillates to the desired station’s broadcast (or driving)
frequency. The more selective the radio is in discriminating between stations, the smaller its damping. A child on a swing is
driven by a parent at the swing’s natural frequency to achieve maximum amplitude. In all of these cases, the efficiency of
energy transfer from the driving force into the oscillator is best at resonance. On gravel roads that are corrugated, if people
travel at the ‘wrong’ speed, the bumps are very noticeable. At other speeds, it is difficult to feel the bumps at all. shows a
photograph of a famous example (the Tacoma Narrows Bridge) of the destructive effects of a driven harmonic oscillation.
Unfortunately, heavy winds happened to drive the bridge at its natural frequency, leading to the collapse.
Collapse of the Tacoma Narrows Bridge: In 1940, the Tacoma Narrows Bridge in Washington state collapsed. Heavy cross
winds drove the bridge into oscillations at its resonant frequency. The damping decreased when support cables broke loose and
started to slip over the towers, allowing increasingly greater amplitudes until the structure failed. (credit: PRI’s Studio 360, via
Flickr)
Key Points
The waves either interfere with each other constructively or destructively, which will either amplify or minimize the wave,
respectively.
Spherical waves are emitted from a single point source in a spherical shape.
A plane wave is a constant-frequency wave whose wavefronts (surfaces of constant phase ) are infinite parallel planes of
constant peak-to-peak amplitude normal to the phase velocity vector.
Although it is not possible in practice to have a true plane wave, many waves approximate plane wave behavior.
The reflected wave is inverted from the incident wave when a transverse wave on a string is fixed at the end point. The reflected
wave is not inverted from the incident wave when a transverse wave on a string is free at the end point.
A standing wave occurs when an incident wave meets a reflected wave on a string.
A standing wave contains nodes (points that remain flat due to the destructive interference ) and antinodes (points with
maximum oscillation due to the constructive interference ).
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Every point in the string oscillates up and down and the amplitude of the oscillations depends on the location of the point.
A standing wave has some points that remain flat due to destructive interference. These are called antinodes.
The points on a standing wave that have reached maximum oscillation do so from constructive interference, and are called
nodes.
A standing wave in an air column is a transverse wave.
A node occurs at the closed end of an air tube where there is no wave displacement.
An antinode occurs at the open end of the air tube where the maximum displacement occurs.
The resonant frequencies of a tube closed at one end are as follows: f = n , n=1,3,5.
n
vw
4L
Any oscillator has a natural frequency. Driven at the natural frequency, oscillations increase in amplitude with each oscillation.
The widths of the resonance curves depend on damping: less damping corresponds to a narrower resonance. With less damping,
the resonance peak becomes more pronounced.
The relationship between resonance and damping applies to a huge variety of systems. Damping is often reduced or enhanced to
induce the desired response of an oscillator.
Key Terms
wavefront: An imaginary surface passing through points of a medium oscillating in phase.
destructive interference: Occurs when waves interfere with each other crest to trough (peak to valley) and are exactly out of
phase with each other.
constructive interference: Occurs when waves interfere with each other crest to crest and the waves are exactly in phase with
each other.
transverse wave: Any wave in which the direction of disturbance is perpendicular to the direction of travel.
natural frequency: The frequency at which a system vibrates on its own. For a spring (spring constant k) with an object of
−
−
mass m attached, the natural frequency is given as fn =
1
2π
√
k
m
.
damping: The reduction in the magnitude of oscillations by the dissipation of energy
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CHAPTER OVERVIEW
This page titled 17: Electric Charge and Field is shared under a not declared license and was authored, remixed, and/or curated by Boundless.
1
17.1: Overview
learning objectives
Identify factors that determine the atom’s net charge
Protons
Protons are found in the center of the atom; they, with neutrons, make up the nucleus. Protons have a charge of +1 and a mass of 1
atomic mass unit, which is approximately equal to 1.66×10-24 grams. The number of protons in an atom defines the identity of the
element (an atom with 1 proton is hydrogen, for example, and an atom with two protons is helium). As such, protons are relatively
stable; their number rarely changes, only in the instance of radioactive decay.
Electrons
Electrons are found in the periphery of the atom and have a charge of -1. They are much smaller than protons; their mass is
1183611836 amu. Typically in modeling atoms, protons and neutrons are regarded as stationary, while electrons move about in the
space outside the nucleus like a cloud. The negatively charged electronic cloud indicates the regions of the space where electrons
are likely to be found. The electrons cloud patterns are extremely complex and is of no importance to the discussion of electric
charge in the atom. More important is the fact that electrons are labile; that is, they can be transferred from one atom to the next. It
is through electronic transfer that atoms become charged.
Ions
In the ground state, an atom will have an equal number of protons and electrons, and thus will have a net charge of 0. However,
because electrons can be transferred from one atom to another, it is possible for atoms to become charged. Atoms in such a state are
known as ions.
If a neutral atom gains an electron, it becomes negative. This kind of ion is called an anion.
If a neutral atom loses an electron, it becomes positive. This kind of ion is called a cation.
The steady flow of electrons is called current. Current is what flows through electrical wires and powers electronics items, from
light bulbs to televisions.
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High School Physics - Electric Charge
Planetary Model of an Atom: Small electrons orbit the large and relatively fixed nucleus of protons and neutrons.
17.1.2 https://phys.libretexts.org/@go/page/14543
Properties of Electric Charges
Electric charge is a fundamental physical property of matter that has many parallels to mass.
learning objectives
Describe properties of electric charge, such as its relativistic invariance and its conservation in closed systems
Invariance
Like mass, electric charge in a closed system is conserved. As long as a system is impermeable, the amount of charge inside it will
neither increase nor decrease; it can only be transferred. However, electric charge differs from other properties—like mass—in that
it is a relativistic invariant. That is, charge is independent of speed. The mass of a particle will rise exponentially as its speed
approaches that of light, its charge, however, will remain constant.
The independence of electric charge from speed was proven through an experiment in which one fast-moving helium nucleus (two
protons and two neutrons bound together) was proven to have the same charge as two separate, slow-moving deuterium nuclei (one
proton and one neutron bound together in each nucleus).
Coulomb’s Law: The forces (F1 and F2) sum to produce the total force, which is calculated by Coulomb’s Law and is proportional
to the product of the charges q1 and q2, and inversely proportional to the square of the distance (r21) between them.
The formula for gravitational force has exactly the same form as Coulomb’s Law, but relates the product of two masses (rather than
the charges) and uses a different constant. Both act in a vacuum and are central (depend only on distance between the forces) and
conservative (independent of path taken). However, it should be noted that when comparing similar terms, charge-based interaction
is substantially greater than that based on mass. For example, the electric repulsion between two electrons is about 1042 times
stronger than their gravitational attraction.
Charge Separation
Charge separation, often referred to as static electricity, is the building of space between particles of opposite charges.
learning objectives
Identify factors that can create charge separation
All matter is composed of atoms made up of negatively-charged electrons and positively-charged protons. In the ground state, each
atom is of neutral charge—its protons and electrons are equal in number, and it exists with no permanent dipole. Because electrons
are labile (i.e., they can be transferred from atom to atom) it is possible for the phenomenon of “charge separation” (often referred
to as static electricity) to occur.
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Static Electricity: Due to friction between her hair and the plastic slide, the girl on the left has created charge separation, resulting
in her hair being attracted to the slide.
In chemistry, this charge separation is illustrated simply by the transfer of an electron from one atom to another as an ionic bond is
formed. In physics, there are many other instances of charge separation that cannot be written as formal chemical reactions.
Consider, for example, rubbing a balloon on your hair. Once you pull the balloon away, your hair will stand on end and “reach”
towards the balloon. This is because electrons from one have transferred to the other, causing one to be positive and the other to be
negative. Thus, the opposite charges attract. A similar example can be seen in playground slides (as shown in ).
Charge separation can be created not only by friction, but by pressure, heat, and other charges. Both pressure and heat increase the
energy of a material and can cause electrons to break free and separate from their nuclei. Charge, meanwhile, can attract electrons
to or repel them from a nucleus. For example, a nearby negative charge can “push” electrons away from the nucleus around which
they typically orbit. Charge separation occurs often in the natural world. It can have an extreme effect if it reaches a critical level,
whereat it becomes discharged. Lightning is a common example.
Polarization
Dielectric polarization is the phenomenon that arises when positive and negative charges in a material are separated.
learning objectives
Identify two ways polarization can occur on the molecular level
The concept of polarity is very broad and can be applied to molecules, light, and electric fields. For the purposes of this atom, we
focus on its meaning in the context of what is known as dielectric polarization—the separation of charges in materials.
Dielectrics
A dielectric is an insulator that can be polarized by an electric field, meaning that it is a material in which charge does not flow
freely, but in the presence of an electric field it can shift its charge distribution. Positive charge in a dielectric will migrate towards
the applied field, while negative charges will shift away. This creates a weak local field within the material that opposes the applied
field.
Different materials will react differently to an induced field, depending on their dielectric constant. This constant is the degree of
their polarizability (the extent to which they become polarized).
Atomic Model
The most basic view of dielectrics involves considering their charged components: protons and electrons. If an electric field is
applied to an atom, the electrons in the atom will migrate away from the applied field. The protons, however, remain relatively
exposed to the field. This separation creates a dipole moment, as shown in.
17.1.4 https://phys.libretexts.org/@go/page/14543
Reaction of an Atom to an Applied Electric Field: When an electric field (E) is applied, electrons drift away from the field. Their
average location is displaced from the average location of the protons (which hasn’t moved) by a distance of d. The atom’s dipole
moment is represented by M.
Dipole Polarization
On the molecular level, polarization can occur with both dipoles and ions. In polar bonds, electrons are more attracted to one
nucleus than to the other. One example of a dipole molecule is water, (H2O), which has a bent shape (the H-O-H angle is 104.45°)
and in which the oxygen pulls electron density away from the H atoms, leaving the H relatively positive and the O relatively
negative, as shown in.
Water Molecule: Water is an example of a dipole molecule, which has a bent shape (the H-O-H angle is 104.45°) and in which the
oxygen pulls electron density away from the H atoms, leaving the H relatively positive and the O relatively negative.
When a dipolar molecule is exposed to an electric field, the molecule will align itself with the field, with the positive end towards
the electric field and the negative end away from it.
Ionic Polarization
Ionic compounds are those that are formed from permanently charge-separated ions. For example, table salt (NaCl) is formed from
Na+ and Cl– ions that are not formally bound to one another through a chemical bond, but interact very strongly due to their
opposite charges.
Ions are still free from one another and will naturally move at random. If they happen to move in a way that is asymmetrical, and
results in a greater concentration of positive ions in one area and a greater concentration of negative ions in another, the sample of
ionic compound will be polarized—a phenomenon is known as ionic polarization.
learning objectives
Formulate rules that apply to the creation and the destruction of electric charge
Electric charge is a physical property of matter. It is created by an imbalance in a substance’s number of protons and electrons. The
matter is positively charged if it contains more protons than electrons, and it is negatively charged if it contains more electrons than
17.1.5 https://phys.libretexts.org/@go/page/14543
protons. In both instances, charged particles will experience a force when in the presence of other charged matter.
Charges of like sign (positive and positive, or negative and negative) will repel each other, whereas charges of opposite sign
(positive and negative) will attract each another, as shown in.
Charge Repulsion and Attraction: Charges of like sign (positive and positive, or negative and negative) will repel each other,
whereas charges of opposite sign (positive and negative) will attract each other.
The SI unit for charge is the Coulomb (C), which is approximately equal to 6.24×10186.24×1018 elementary charges. (An
elementary charge is the magnitude of charge of a proton or electron. )
Conservation of Charge
Charge, like matter, is essentially constant throughout the universe and over time. In physics, charge conservation is the principle
that electric charge can neither be created nor destroyed. The net quantity of electric charge, the amount of positive charge minus
the amount of negative charge in the universe, is always conserved.
For any finite volume, the law of conservation of charge (Q) can be written as a continuity equation:
Q (t2 ) = Q (t1 ) + Q in − Q out (17.1.2)
where Q(t1) is the charge in the system at a given time, Q(t2) is the charge in the same system at a later time, Qin is the charge that
has entered the system between the two times, and Qout is the amount of charge that has left the system between the two times.
This does not mean that individual positive and negative charges cannot be created or destroyed. Electric charge is carried by
subatomic particles such as electrons and protons, which can be created and destroyed. For example, when particles are destroyed,
equal numbers of positive and negative charges are destroyed, keeping the net amount of charge unchanged.
Static Electricity
Static electricity is when an excess of electric charge collects on an object’s surface. It can be created through contact between
materials, a buildup of pressure or heat, or the presence of a charge. Static electricity can also be created through friction between a
balloon (or another object) and human hair (see ). It can be observed in storm clouds as a result of pressure buildup; lightning (see )
is the discharge that occurs after the charge exceeds a critical concentration.
Static Electricity: Due to friction between her hair and the plastic slide, the girl on the left has created charge separation, resulting
in her hair being attracted to the slide.
17.1.6 https://phys.libretexts.org/@go/page/14543
Lightning: Lightning is a dramatic natural example of static discharge.
learning objectives
Identify conductors and insulators among common materials
Overview
All materials can be categorized as either insulators or conductors based on a physical property known as resistivity.
An insulator is a material in which, when exposed to an electric field, the electric charges do not flow freely—it has a high
resistivity. Conversely, a conductor is a material that permits the flow of electric charges in one or more directions—its resistivity is
low.
Conductors
All conductors contain electric charges that, when exposed to a potential difference, move towards one pole or the other. The
positive charges in a conductor will migrate towards the negative end of the potential difference; the negative charges in the
material will move towards the positive end of the potential difference. This flow of charge is electric current.
Ionic substances and solutions can conduct electricity, but the most common and effective conductors are metals. Copper is
commonly used in wires due to its high conductivity and relatively inexpensive price. However, gold-plated wires are sometimes
used in instances in which especially high conductivity is necessary.
Every conductor has a limit to its ampacity, or amount of current it can carry. This usually is the current at which the heat released
due to resistance melts the material.
Insulators
Insulators are materials in which the internal charge cannot flow freely, and thus cannot conduct electric current to an appreciable
degree when exposed to an electric field.
While there is no perfect insulator with infinite resistivity, materials like glass, paper and Teflon have very high resistivity and can
effectively serve as insulators in most instances.
Just as conductors are used to carry electrical current through wires, insulators are commonly used as coating for the wires.
Insulators, like conductors, have their physical limits. When exposed to enough voltage, an insulator will experience what is known
as electrical breakdown, in which current suddenly spikes through the material as it becomes a conductor.
17.1.7 https://phys.libretexts.org/@go/page/14543
Conductor and Insulator in a Wire: This wire consists of a core of copper (a conductor) and a coating of polyethylene (an
insulator). The copper allows current to flow through the wire, while the polyethylene ensures that the current does not escape.
learning objectives
Explain the difference in value of a real electron’s charge and the charge measured by Robert Millikan
Simplified scheme of Millikan’s oil-drop experiment: This apparatus has a parallel pair of horizontal metal plates. A uniform
electric field is created between them. The ring has three holes for illumination and one for viewing through a microscope. Special
oil for vacuum apparatus is sprayed into the chamber, where drops become electrically charged. The droplets enter the space
between the plates and can be controlled by changing the voltage across the plates.
Millikan then calculated the charge on particles suspended in mid-air. His assumptions were that the force of gravity, which is the
product of mass (m) and gravitational acceleration (g), was equal to the force of the electric field (the product of the charge (q) and
the electric field (E)):
q⋅E = m⋅g (17.1.3)
m⋅g
q = (17.1.4)
E
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Since he already knew the mass of the oil droplets and the acceleration due to gravity (9.81 m/s^2), as well as the energy of the x-
rays he was using, he was able to calculate the charge.
Although the charge of each droplet was unknown, Millikan adjusted the strength of the X-rays ionizing the air and measured many
values of (q) from many different oil droplets. In each instance, the charge measured was a multiple of 1.5924(17)×10−19 C. Thus,
it was concluded that the elementary electric charge was 1.5924(17)×10−19 C.
The results were very accurate. The calculated value from the Oil-Drop Experiment differs by less than one percent of the current
accepted value of 1.602176487(40)×10−19 C.
The Oil-Drop Experiment was tremendously influential at the time, not only for determining the charge of an electron, but for
helping prove the existence of particles smaller than atoms. At the time, it was not fully accepted that protons, neutrons, and
electrons existed.
Key Points
A proton is a positively charged particle located in the nucleus of an atom. An electron has 1
1836
times the mass of a proton, but
an equal and opposite negative charge.
An elementary charge — that of a proton or electron — is approximately equal to 1.6×10-19Coulombs.
Unlike protons, electrons can move from atom to atom. If an atom has an equal number of protons and electrons, its net charge
is 0. If it gains an extra electron, it becomes negatively charged and is known as an anion. If it loses an electron, it becomes
positively charged and is known as a cation.
Charge is measured in Coulombs (C), which represent 6.242×1018 e, where e is the charge of a proton. Charges can be positive
or negative, and as such a singular proton has a charge of 1.602×10−19 C, while an electron has a charge of -1.602×10−19 C.
Electric charge, like mass, is conserved. The force generated by two charges is of the same form as that generated by two
masses and, like gravity, force from an electrical field is both conservative and central.
Electric charge is a relativistic invariant. That is, charge (unlike mass) is independent of speed. Whereas the mass of a particle
will exponentially rise as its speed approaches that of light, charge will remain constant.
Because electrons are labile (i.e., they can be transferred from atom to atom), it is possible for “charge separation” to occur.
This phenomenon is often commonly referred to as static electricity.
Charge separation can be created by friction, pressure, heat, and other charges.
Charge separation can reach a critical level, whereat it is discharged. Lightning is a common example.
Dielectrics are insulators that are capable of being polarized by an electric field. That is, their charges cannot flow freely, but
can still be induced to redistribute unevenly.
Electric fields applied to atoms will push electrons away from the field. In the case of polar molecules, the negative ends thereof
will align themselves away from the field while the positive ends will be towards the field.
An instantaneous polarization occurs when ions, through natural, random vibrations, become distributed asymmetrically such
that one area is more dense with one type of ion than another.
Electric charge is a physical property of matter created by an imbalance in the number of protons and electrons in a substance.
Charge can be created or destroyed. However, any creation or elimination of charge occurs at a ratio of 1:1 between positive
and negative charges.
Static electricity is when an excess of electric charge collects on an object’s surface.
Resistivity, a physical property that measures the ability of a material to carry current, is the main factor in determining whether
a substance is a conductor or an insulator.
Conductors contain electric charges that, when exposed to a potential difference, move towards one pole or the other. This flow
of charge is electric current.
Insulators are materials in which the internal charge cannot flow freely, and thus cannot conduct electric current to an
appreciable degree when exposed to an electric field.
The Oil-Drop Experiment involved ionizing droplets of oil as they fell through the air, and balancing the force of gravity with
the force of an electric field applied by electrodes above and below the droplet.
Millikan could not directly count the number of electrons on each oil droplet, but found that the common denominator between
all measured charges was equal to 1.5924(17)×10−19 C, and thus concluded that this value was the charge of an electron.
The measured value of an electron’s charge, 1.5924(17)×10−19 C, differs from the accepted value of 1.602176487(40)×10−19 C
by less than one percent.
17.1.9 https://phys.libretexts.org/@go/page/14543
Key Terms
nucleus: the massive, positively charged central part of an atom, made up of protons and neutrons
coulomb: In the International System of Units, the derived unit of electric charge; the amount of electric charge carried by a
current of 1 ampere flowing for 1 second. Symbol: C
gravity: Resultant force on Earth’s surface, of the attraction by the Earth’s masses, and the centrifugal pseudo-force caused by
the Earth’s rotation.
electric field: A region of space around a charged particle, or between two voltages; it exerts a force on charged objects in its
vicinity.
discharge: the act of releasing an accumulated charge
static electricity: an electric charge that has built up on an insulated body, often due to friction
dipole moment: The vector product of the charge on either pole of a dipole and the distance separating them.
dielectric: An electrically insulating or nonconducting material considered for its electric susceptibility (i.e., its property of
polarization when exposed to an external electric field).
insulator: A substance that does not transmit heat (thermal insulator), sound (acoustic insulator) or electricity (electrical
insulator).
electric charge: A quantum number that determines the electromagnetic interactions of some subatomic particles; by
convention, the electron has an electric charge of -1 and the proton +1, and quarks have fractional charge.
conductor: A material which contains movable electric charges.
resistivity: In general, the resistance to electric current of a material; in particular, the degree to which a material resists the
flow of electricity.
voltage: The amount of electrostatic potential between two points in space.
terminal velocity: The speed at which an object in free-fall and not in a vacuum ceases to accelerate downwards because the
force of gravity is equal and opposite to the drag force acting against it.
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17.2: Shelding and Charging Through Induction
learning objectives
Describe structure of a Faraday cage
Electrostatic shielding is the phenomenon that is observed when a Faraday cage operates to block the effects of an electric field.
Such a cage can block the effects of an external field on its internal contents, or the effects of an internal field on the outside
environment.
A Faraday cage is a closed chamber consisting of a conducting material or a mesh of such a material. This type of cage was first
invented by Michael Faraday in 1836, and can block external static and non-static electric fields.
When an external electric field operates on a Faraday cage, the charges within the cage (which are mobile, as the cage is a
conductor) rearrange themselves to directly counteract the field and thus “shield” the interior of the cage from the external field
Faraday Cage in Presence of an External Electrical Field: As the field is applied, the negative charge from the cage migrates
toward the positive end of the field, canceling the effects of the field at both ends of the cage.
The action of a Faraday cage may depend on whether or not it is grounded. Consider a charge placed within a cage. If the cage is
not grounded, electrons in the cage will redistribute such that the interior wall of the cage takes on a charge opposite the internal
charge. This would leave an exterior wall of opposite charge to that of the interior. If it is grounded, however, excess charges on the
exterior of the cage will go to the ground, leaving the exterior wall of neutral charge.
Limitations
Faraday cages are limited in their effectiveness, and cannot block static and slowly varying magnetic fields, such as that of the
planet Earth. They can, however, shield the interior from external magnetic radiation provided that the mesh is smaller than the
wavelength of the radiation and that the shield is sufficiently thick.
Applications
Microwave ovens contain energy within themselves, shielding the outside from harmful radiation.
Electrical linemen often wear suits made of Faraday cages so as to avoid electrocution.
Elevators can act as unintended Faraday cages, shielding cell phones and radios from signal from the outside.
Induced Charge
Electrostatic induction is the redistribution of charges within an object that occurs as a reaction to the presence of a nearby charge.
learning objectives
Compare electrostatic induction processes in conductors and dielectrics
Electrostatic induction is the redistribution of charge within an object, which occurs as a reaction to a nearby charge.
Normally, a unit of matter will have equal parts positive and negative charge, distributed evenly throughout the object. As such, it
has no net charge.
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If a charged piece of matter is brought into close proximity with respect to an uncharged piece of matter, it can induce charge
redistribution in the neutral material.
In such an event, the electrons in the neutral object move (the protons are relatively inert) according to the charge of the nearby
charged object (inductor). If the inductor is positive, electrons migrate toward it, making the uncharged object more negative in that
area and positive in the region opposite it. If the inductor is negative, the electrons in the neutral object are repelled, leaving a
positive charge near the inductor and a negative charge opposite it.
Electric Induction Experiment: Circa 1870, the positive end of an electrostatic generator is placed near an uncharged brass
cylinder, causing the cylinder to polarize as its left end becomes positive and its right end becomes negative.
If a charged object comes in touch with an uncharged object, or the two come sufficiently close to one another to cause a discharge
that bridges the gap between them, the previously uncharged object will become charged. Depending on the sign of the charge of
the inductor, electrons will go to or leave the previously uncharged object. Total charge is conserved, and that of the inductor
decreases as it transfers charge to its subject.
Subjects that can react to inductors include conductors and dielectrics. In the case of the former, the free flow of charges makes it
possible for strong polarization to occur. In the case of the latter, the force is comparatively weak.
Key Points
Electrostatic shielding is the phenomenon that is observed when a Faraday cage operates to block the effects of an electric field.
Such a cage can block the effects of an external field on its internal contents, or the effects of an internal field on the outside
environment.
A Faraday cage is a closed chamber consisting of a conducting material or a mesh of such a material. This type of cage can
block external static and non-static electric fields.
Faraday cages cannot block static and slowly varying magnetic fields, such as that of the planet Earth. They can, however,
shield the interior from external magnetic radiation provided that the mesh is smaller than the wavelength of the radiation and
that the shield is sufficiently thick.
If a charged piece of matter is brought into close proximity with respect to an uncharged piece of matter, it can induce charge
redistribution in the neutral material. This is one form of induction.
If a charged object comes in touch with an uncharged object, or the two come sufficiently close to one another to cause a
discharge that bridges the gap between them, the previously uncharged object will become charged. This is another form of
induction.
Subjects that can react to inductors include conductors and dielectrics. In the case of the former, the free flow of charges makes
it possible for strong polarization to occur. In the case of the latter, the force is comparatively weak.
Key Terms
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ , and corresponds to the velocity of the wave divided by its frequency.
discharge: the act of releasing an accumulated charge
inductor: A passive device that introduces inductance into an electrical circuit.
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dielectric: An electrically insulating or nonconducting material considered for its electric susceptibility (i.e., its property of
polarization when exposed to an external electric field).
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17.3: Coulomb’s Law
learning objectives
Apply the superposition principle to determine the net response caused by two or more stimuli
The superposition principle (also known as superposition property) states that: for all linear systems, the net response at a given
place and time caused by two or more stimuli is the sum of the responses which would have been caused by each stimulus
individually. For Coulomb’s law, the stimuli are forces. Therefore, the principle suggests that total force is a vector sum of
individual forces.
Coulomb Force
The scalar form of Coulomb’s Law relates the magnitude and sign of the electrostatic force F, acting simultaneously on two point
charges q1 and q2:
1 | q1 q2 |
|F| = (17.3.1)
2
4πarϵ0 r
Lorentz Force on a Moving Particle: Lorentz force f on a charged particle (of charge q) in motion (instantaneous velocity v). The
E field and B field vary in space and time.
where r is the separation distance and ε0 is electric permittivity. If the product q1q2 is positive, the force between them is repulsive;
if q1q2 is negative, the force between them is attractive. The principle of linear superposition allows the extension of Coulomb’s
law to include any number of point charges—in order to derive the force on any one point charge by a vector addition of these
individual forces acting alone on that point charge. The resulting force vector happens to be parallel to the electric field vector at
that point, with that point charge removed.
To calculate the force on a small test charge q at position rr, due to a system of N discrete charges:
N N
^
q r − ri q Ri
F(r) = ∑ qi = ∑ qi (17.3.2)
3 2
4πarϵ0 |r − ri | r4πarϵ0 | Ri |
i=1 i=1
where qi and ri are the magnitude and position vector of the i-th charge, respectively, and ^
Ri is a unit vector in the direction of
R = r−r
i i (a vector pointing from charges qi to q. )
Of course, our discussion of superposition of forces applies to any types (or combinations) of forces. For example, when a charge is
moving in the presence of a magnetic field as well as an electric field, the charge will feel both electrostatic and magnetic forces.
Total force, affecting the motion of the charge, will be the vector sum of the two forces. (In this particular example of the moving
charge, the force due to the presence of electromagnetic field is collectively called Lorentz force (see ).
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learning objectives
Describe shape of a Coulomb force from a spherical distribution of charge
Through the work of scientists in the late 18th century, the main features of the electrostatic force —the existence of two types of
charge, the observation that like charges repel, unlike charges attract, and the decrease of force with distance—were eventually
refined, and expressed as a mathematical formula. The mathematical formula for the electrostatic force is called Coulomb ‘s law
after the French physicist Charles Coulomb (1736–1806), who performed experiments and first proposed a formula to calculate it.
Charge distribution in a water molecule: Schematic representation of the outer electron cloud of a neutral water molecule. The
electrons spend more time near the oxygen than the hydrogens, giving a permanent charge separation as shown. Water is thus a
polar molecule. It is more easily affected by electrostatic forces than molecules with uniform charge distributions.
Modern experiments have verified Coulomb’s law to great precision. For example, it has been shown that the force is inversely
proportional to distance between two objects squared (F∝1/r2) to an accuracy of 1 part in 1016. No exceptions have ever been
found, even at the small distances within the atom.
Coulomb’s law holds even within the atoms, correctly describing the force between the positively charged nucleus and each of the
negatively charged electrons. This simple law also correctly accounts for the forces that bind atoms together to form molecules and
for the forces that bind atoms and molecules together to form solids and liquids.
Generally, as the distance between ions increases, the energy of attraction approaches zero and ionic bonding is less favorable. As
the magnitude of opposing charges increases, energy increases and ionic bonding is more favorable.
An electric field is a vector field which associates to each point of the space the Coulomb force that will experience a test unity
charge. Given the electric field, the strength and direction of a force F on a quantity charge q in an electric field E is determined by
the electric field. For a positive charge, the direction of the electric field points along lines directed radially away from the location
of the point charge, while the direction is towards for a negative charge.
This distribution around a charged molecule is spherical in nature, and creates a sort of electrostatic “cloud” around the molecule.
The attraction or repulsion forces within the spherical distribution of charge is stronger closer to the molecule, and becomes weaker
as the distance from the molecule increases.
This image shows the outer electron cloud of a neutral water molecule. The charge distribution of the oxygen molecule is negative,
and attracts the two positive hydrogen molecules. The attraction between the two opposing charges forms a neutral water molecule.
It is a polar molecule because there is still a permanent charge separation because the electrons spend more time near the oxygen
than the hydrogens.
learning objectives
Explain when the vector notation of Coulomb’s Law can be used
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Solving Problems with Vectors and Coulomb’s Law
Electric Force Between Two Point Charges
To address the electrostatic forces among electrically charged particles, first consider two particles with electric charges q and Q,
separated in empty space by a distance r. Suppose that we want to find the electric force vector on charge q. (The electric force
vector has both a magnitude and a direction. ) We can express the location of charge q as rq, and the location of charge Q as rQ. In
this way we can know both how strong the electric force is on a charge, but also what direction that force is directed in. Coulomb’s
Law using vectors can be written as:
kq Q (rq − rQ )
FE = (17.3.3)
3
| rq − rQ |
which the charges are immersed. The variables F , Γ and R are in bold because they are vectors. Thus, we need to find
E q Q
r −r
q Q by performing standard vector subtraction. This means that we need to subtract the corresponding components of vector
\mathbf{r}_\mathrm{Q}} from vector \mathrmbfr . q
This vector notation can be used in the simple example of two point charges where only one of which is a source of charge.
Application of Coulomb’s Law: In a simple example, the vector notation of Coulomb’s Law can be used when there are two point
charges and only one of which is a source charge.
Multiple point charges: Coulomb’s Law applied to more than one point source charges providing forces on a field charge.
Note the coordinate system that has been chosen. All of the charges lie on the corners of a square, and the origin is chosen to
collocate with the lower right source charge, and aligned with the square. Since we can have only one origin of coordinates, no
more than one of the source points can lie at the origin, and the displacements from different source points to the field point differ.
The total force on the field charge q is due to applications of the force described in the vector notation of Coulomb’s Law from
each of the source charges. The total force is therefore the sum of these individual forces.
Displacements of field charge: The displacements of the field charge from each source charge are shown as light blue arrows.
Applying Coulomb’s Law three times and summing the results gives us:
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kq ⋅ q1 (rq − rq1 ) kq ⋅ q2 (rq − rq2 ) kq ⋅ q3 (rq − rq3 )
FEq = + + (17.3.4)
3 3 3
| rq − rq1 | | rq − rq2 | | rq − rq3 |
This equation can further be simplified and applied to a fixed number of charge points.
qn qi (rn − ri )
Fn = ∑ (17.3.5)
3
i≠n 4πϵ0 | rn − ri |
Coulomb’s Law: In this video I continue with my series of tutorial videos on Electrostatics. It’s pitched at undergraduate level and
while it is mainly aimed at physics majors, it should be useful to anybody taking a first course in electricity and magnetism such as
engineers etc.
Key Points
The superposition principle suggests that the net response at a given place and time caused by two or more stimuli is the sum of
the responses which would have been caused by each stimulus individually.
Total Coulomb force on a test charge due to a group of charges is equal to the vector sum of all the Coulomb forces between the
test charge and other individual charges.
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The superposition of forces is not limited to Coulomb forces. It applies to any types (or combinations) of forces.
The force between two objects is inversely proportional to the square of the distance between two objects.
The attraction or repulsion forces within the spherical distribution of charge is stronger closer to the molecule and becomes
weaker as the distance from the molecule increases.
This law also accounts for the forces that bind atoms together to form molecules and for the forces that bind atoms and
molecules together to form solids and liquids.
The vector notation of Coulomb ‘s Law can be used in the simple example of two point charges where only one of which is a
source of charge.
The total force on the field charge for multiple point source charges is the sum of these individual forces.
Coulomb’s Law can be further simplified and applied to a fixed number of charge points.
Key Items
Lorentz force: The force exerted on a charged particle in an electromagnetic field.
unit vector: A vector with length 1.
electrostatic force: The electrostatic interaction between electrically charged particles; the amount and direction of attraction or
repulsion between two charged bodies.
coulomb’s law: the mathematical equation calculating the electrostatic force vector between two charged particles
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17.4: The Electric Field Revisited
learning objectives
Identify law that can be used to calculate an electric field created by a point charge
The electric field of a point charge is, like any electric field, a vector field that represents the effect that the point charge has on
other charges around it. The effect is felt as a force, and when charged particles are not in motion, this force is known as the
electrostatic force. The electrostatic force is, much like gravity, a force that acts at a distance. Therefore, we rationalize this action
at a distance by saying that charges create fields around them that have effects on other charges.
Given a point charge, or a particle of infinitesimal size that contains a certain charge, electric field lines emanate radially in all
directions. If the charge is positive, field lines point radially away from it; if the charge is negative, field lines point radially
towards it.
Electric field of positive point charge: The electric field of a positively charged particle points radially away from the charge.
Electric field of negative point charge: The electric field of a negatively charged particle points radially toward the particle.
The reason for these directions can be seen in the derivation of the electric field of a point charge. Let’s first take a look at the
definition of the electric field of a point particle:
1 q q
⃗ ^ ^
E = r =k r (17.4.1)
4πϵo r2 2
r
The above equation is defined in radial coordinates, which can be seen in. The constant k is a result of simply combining the
constants together, and q is the charge of the particle creating the electric field. This charge is either positive or negative. If the
17.4.1 https://phys.libretexts.org/@go/page/14546
charge is positive, as shown above, the electric field will be pointing in a positive radial direction from the charge q (away from the
charge).
As a demonstration of this phenomenon, if we now place another positive charge, Q (called the test charge), at some radial
distance, R, away from the original particle, the test charge will feel a force given by
1 q
⃗ ⃗ ^
F = QE = Q r (17.4.2)
4πϵo R2
Radial Coordinate System: The electric field of a point charge is defined in radial coordinates. The positive r direction points
away from the origin, and the negative r direction points toward the origin. The electric field of a point charge is symmetric with
respect to the θ direction.
The thing to keep in mind is that the force above is acting on the test charge Q, in the positive radial direction as defined by the
original charge q. This means that because the charges are both positive and will repel one another, the force on the test charge
points away from the original charge.
If the test charge were negative, the force felt on that charge would be:
1 q
⃗ ⃗ ^
F = QE = −Q r (17.4.3)
4πϵ0 R2
Notice that this points in the negative ^rr^ direction, which is toward the original charge. This makes sense because opposite
charges attract, and the force on the test charge will tend to push it toward the original positive charge creating the field. The above
mathematical description of the electric field of a point charge is known as Coulomb ‘s law.
Superposition of Fields
The resultant of multiple electric fields acting on the same point is the sum of the strength of each field’s applied force at that point.
learning objectives
Formulate the superposition principle for a linear system
As vector fields, electric fields obey the superposition principle. This principle states that for all linear systems, the net response to
multiple stimuli at a given place and time is equal to the sum of the responses that would have resulted from each stimulus
individually.
Possible stimuli include but are not limited to: numbers, functions, vectors, vector fields, and time-varying signals. It should be
noted that the superposition principle is applicable to any linear system, including algebraic equations, linear differential equations,
and systems of equations of the aforementioned forms.
For example, if forces A and B are constant and simultaneously act upon an object, illustrated as O in, the resultant force will be
the sum of forces A and B. Vector addition is commutative, so whether adding A to B or B to A makes no difference on the
resultant vector; this is also the case for subtraction of vectors.
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Vector addition: Forces a and b act upon an object at point O. Their sum is commutative, and results in a resultant vector c.
Electric fields are continuous fields of vectors, so at a given point, one can find the forces that several fields will apply to a test
charge and add them to find the resultant. To do this, first find the component vectors of force applied by each field in each of the
orthogonal axes. This can be done using trigonometric functions. Then once the component vectors are found, add the components
in each axis that are applied by the combined electric fields.
This is one only form of solution. An overall resultant vector can be found by using the Pythagorean theorem to find the resultant
(the hypotenuse of the triangle created with applied forces as legs) and the angle with respect to a given axis by equating the
inverse tangent of the angle to the ratio of the forces of the adjacent and opposite legs.
learning objectives
Calculate the resultant force of the multiple electric charges on a test charge
Thus far, we have looked at electric field lines pertaining to isolated point charges. But what if another charge is introduced? Each
will have its own electric field, and the two fields will interact.
When modeling the electric fields of multiple charges, it’s important to take into consideration the sign and magnitude of each
charge. Such models are not meant to be absolute, but must be self-consistent. For example, the number field of lines should be
proportional to the value of the charge that gives rise to them. This means that if charges q1 (with a +1 value) q2 (+2 charge) and q3
(+3 charge) are in the same field, one can connect 4, 8, and 12 field lines, respectively, to the charges. One could also choose to
connect 3, 6, and 9 field lines, respectively, to q1, q2, and q3; what matters is that the number of lines are related to the charge
values by the same proportionality constant. Field lines should always point away from positive charges and towards negative
charge.
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Field lines between like and unlike charges: Example a shows how the electric field is weak between like charges (the
concentration of field lines is low between them). Example b, by contrast, has a strong field between the charges, as exhibited by
the high concentration of field lines connecting them.
If there are opposite charges in consideration, connect one to the other with field lines. If charges are the same, do not connect them
in any way.
The strength of the electric field depends proportionally upon the separation of the field lines. More field lines per unit area
perpendicular to the lines means a stronger field. It should also be noted that at any point, the direction of the electric field will be
tangent to the field line.
Parallel-Plate Capacitor
A parallel-plate capacitor is an electrical component used to store energy in an electric field between two charged, flat surfaces.
17.4.4 https://phys.libretexts.org/@go/page/14546
learning objectives
Describe general structure of a capacitor
Overview
A capacitor is an electrical component used to store energy in an electric field. Capacitors can take many forms, but all involve two
conductors separated by a dielectric material. For the purpose of this atom, we will focus on parallel-plate capacitors.
Diagram of a Parallel-Plate Capacitor: Charges in the dielectric material line up to oppose the charges of each plate of the
capacitor. An electric field is created between the plates of the capacitor as charge builds on each plate.
Capacitance
All capacitors collect charge on the two, separate conductive surfaces; one side is positive and the other negative. An electric field
is created as charge builds on the opposite surfaces, storing energy. The dielectric between the conductors is meant to act as an
insulator, preventing charge from bridging the gap between the two plates. Such dielectrics are commonly composed of glass, air,
paper, or empty space (a vacuum). In practice, dielectrics do not act as perfect insulators, and permit a small amount of leakage
current to pass through them.
Capacitors are limited in their ability to prevent charge flow from one conductive surface to the other; their ability to hold charge is
measured in Farads (F), which are defined as 1 ampere-second per volt, one joule per square volt and one Coulomb per volt, among
other ways.
For a parallel-plate capacitor, capacitance (C) is related to dielectric permittivity (ε), surface area (A), and separation between the
plates (d):
ϵA
C = (17.4.4)
d
Voltage (V) of a capacitor is related to distance between the plates, dielectric permittivity, conductor surface area, and charge (Q)
on the plates:
Qd
V = (17.4.5)
ϵA
Depending on the dielectric strength (Eds) and distance (d) between plates, a capacitor will “break” at a certain voltage (Vbd). This
is calculated by:
Vbd = Eds d (17.4.6)
17.4.5 https://phys.libretexts.org/@go/page/14546
High School Physics - Parallel Plates and Equipotent…
Equipotent…
Parallel Plates and Equipotential Lines: A brief overview of parallel plates and equipotential lines from the viewpoint of
electrostatics.
learning objectives
Describe unique properties expressed by electric fields in the presence of conductors
An electric field, like other fields (e.g., gravitational or magnetic), is a vector field that surrounds an object. Electric fields are
found around electric charges and help determine the direction and magnitude of force the charge exerts on a nearby charged
particle. It measures units of force exerted per unit of charge, and its SI units are N/C.
17.4.6 https://phys.libretexts.org/@go/page/14546
Field Lines Created by a Point Charge: Lines around the positive charge represent the electric field it creates.
Electrical conductors are materials in which internal charges can move freely. Therefore, they can facilitate the flow of charge, or
current. When a conductor is placed in the presence of an electric field, it exhibits some interesting properties:
1. There is no electric field inside a charged conductor. A charged conductor at electrostatic equilibrium will contain charges only
on its outer surface and will have no net electric field within itself. This is because all the charges in such a conductor will
symmetrically oppose other charges within the conductor, causing the net result to sum to 0.
2. Charged surfaces align themselves perpendicularly relative to electric fields. Provided a conductor is at electrostatic
equilibrium, the electric field upon the surface will be aligned perpendicularly with respect to that surface. If there were a
nonzero parallel component of the electric field, with respect to any charge on the surface of a conductor, that charge would
exert a force and would move. If the conductor is at equilibrium, such a force cannot exist, and therefore the direction of the
electric field must be completely perpendicular to the surface.
3. Curvature on the surface of a conductor allows for increased charge concentration. Charge will not necessarily distribute itself
evenly over the surface of a conductor. If the surface of the conductor is flat, charge will be very evenly distributed. But as the
surface becomes more sharply curved, charge can be found more densely packed in areas, even if the conductor is at
electrostatic equilibrium. Charges on a curved surface repel one another less strongly than they would on a smooth surface. This
is because, based on how the charges are positioned, much of the repulsion the charges exert is in the direction away from the
surface of the conductor, not along its surface. And it is harder for charges to be pushed off a surface than along it. Therefore,
the repulsion between charges on a curved surface is weaker.
Electrical Charge at a Sharp Point of a Conductor: Repulsive forces towards the more sharply curved surface on the right aim
more outward than along the surface of the conductor.
learning objectives
Describe behavior of charges in a conductor in the presence of charge or an electric field and under static equilibrium
Conductors are materials in which charges can move freely. If conductors are exposed to charge or an electric field, their internal
charges will rearrange rapidly. For example, if a neutral conductor comes into contact with a rod containing a negative charge,
some of that negative charge will transfer to the conductor at the point of contact. But the charge will not stay local to the contact
point — it will distribute itself evenly over the surface of the conductor. Once the charges are redistributed, the conductor is in a
17.4.7 https://phys.libretexts.org/@go/page/14546
state of electrostatic equilibrium. It should be noted that the distribution of charges depends on the shape of the conductor and that
static equilibrium may not necessarily involve an even distribution of charges, which tend to aggregate in higher concentrations
around sharp points. This is explained in.
Electrical Charge at a Sharp Point of a Conductor: Forces between like charges at either end of the conductor are identical, but
the components of the forces parallel to the surfaces are different. The component parallel to the surface is greatest on the flattest
surface and therefore moves charges away from one another more freely. This explains the difference in concentration of charge on
flat vs. pointed areas of a conductor.
Similarly, if a conductor is placed in an electric field, the charges within the conductor will move until the field is perpendicular to
the surface of the conductor. Negative charges in the conductor will align themselves towards the positive end of the electric field,
leaving positive charges at the negative end of the field. The conductor thus becomes polarized, with the electric field becoming
stronger near the conductor but disintegrating inside it. This occurrence is similar to that observed in a Faraday cage, which is an
enclosure made of a conducting material that shields the inside from an external electric charge or field or shields the outside from
an internal electric charge or field.
Key Points
The electric field is a vector field around a charged particle. It represents the force that other charged particles would feel if
placed near the particle creating the electric field.
Given a point charge, or a particle of infinitesimal size, that contains a certain charge, electric field lines emanate from equally
in all radial directions.
If the point charge is positive, field lines point away from it; if the charge is negative, field lines point toward it.
The superposition principle states that for all linear systems, the net response to multiple stimuli at a given place and time is
equal to the sum of the responses that would have been caused by each stimulus individually.
Possible stimuli include but are not limited to numbers, functions, vectors, vector fields, and time-varying signals.
The superposition principle is applicable to any linear system, including algebraic equations, linear differential equations, and
systems of equations of the aforementioned forms.
Electric fields are continuous fields of vectors, so at a given point, one can find the forces that several fields will apply to a test
charge and add them to find the resultant.
When multiple electric charges interact, their resultant force on a test charge can be calculated using vector addition.
If there are opposite charges in consideration, connect one to the other with field lines. If charges are the same, do not connect
them in any way.
When modeling the electric fields of multiple charges, take into consideration the sign and magnitude of each charge. The
number of field lines should be proportional to the value of the charge that gives rise to them.
Capacitors can take many forms, but all involve two conductors separated by a dielectric material.
All capacitors collect charge on the two, separate conductive surfaces; one side is positive and the other negative. An electric
field is created as charge builds on the opposite surfaces, storing energy. The dielectric acts as an insulator, isolating the charged
surfaces.
The ability of capacitors to hold charge is measured in Farads (F). Capacitors will typically allow a small leakage of current
through the dielectric, but after a certain voltage, the entire capacitor breaks down as the dielectric becomes a conductor.
There is no electric field inside a charged conductor. This is because charges, which are located on the surface of the conductor,
symmetrically oppose one another and sum to 0 in all locations.
Charged surfaces align themselves perpendicularly relative to electric fields to achieve electrostatic equilibrium. If charges are
not distributed as such, they will exert net force upon one another, which will move them. In such an instance, the charges will
not be at static equilibrium.
17.4.8 https://phys.libretexts.org/@go/page/14546
Curvature on the surface of a field allows for increased charge concentration. Much of the repulsion charges exert is in the
direction moving away from the surface of the conductor, rather than along its surface. Charges thus push one another more
weakly along the surface of a curved conductor.
The presence of charge or an electric field will force charges in a conductor to redistribute along the conductor’s surface until
static equilibrium is achieved.
At static equilibrium, charge will be more concentrated in sharp, pointy areas of conductors than elsewhere.
At static equilibrium, the inside of a conductor will be entirely shielded from an external electric field.
Key Terms
coulomb’s law: the mathematical equation calculating the electrostatic force vector between two charged particles
vector field: a construction in which each point in a Euclidean space is associated with a vector; a function whose range is a
vector space
orthogonal: Of two objects, at right angles; perpendicular to each other.
superposition principle: The principle that a linear combination of two or more solutions of an equation is itself a solution; it is
a feature of many physical laws.
vector: A directed quantity, one with both magnitude and direction; the between two points.
capacitor: An electronic component capable of storing an electric charge, especially one consisting of two conductors separated
by a dielectric.
dielectric: An electrically insulating or nonconducting material considered for its electric susceptibility (i.e., its property of
polarization when exposed to an external electric field).
conductor: A material which contains movable electric charges.
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
static equilibrium: the physical state in which all components of a system are at rest and the net force is equal to zero
throughout the system
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17.5: Electric Flux and Gauss’s Law
learning objectives
Express the electric flux for uniform and non-uniform electric fields
Electric flux is the rate of flow of the electric field through a given area (see ). Electric flux is proportional to the number of electric
field lines going through a virtual surface.
Electric Flux: Electric flux visualized. The ring shows the surface boundaries. The red arrows for the electric field lines.
If the electric field is uniform, the electric flux passing through a surface of vector area S is Φ = E ⋅ S = ES cos θ where E is the
E
magnitude of the electric field (having units of V/m), S is the area of the surface, and θ is the angle between the electric field lines
and the normal ( perpendicular ) to S.
For a non-uniform electric field, the electric flux dΦ through a small surface area dS is given by
E dΦE = E ⋅ dS (the electric
field, E, multiplied by the component of area perpendicular to the field).
Gauss’ Law describes the electric flux over a surface S as the surface integral: Φ = ∬ E ⋅ dS where E is the electric field and dS
E
S
is a differential area on the closed surface S with an outward facing surface normal defining its direction.
It is important to note that while the electric flux is not affected by charges that are not within the closed surface, the net electric
field, E, in the Gauss’ Law equation, can be affected by charges that lie outside the closed surface. While Gauss’ Law holds for all
situations, it is only useful for “by hand” calculations when high degrees of symmetry exist in the electric field. Examples include
spherical and cylindrical symmetry.
Electric flux has SI units of volt metres (V m), or, equivalently, newton metres squared per coulomb (N m2 C−1). Thus, the SI base
units of electric flux are kg·m3·s−3·A−1.
Gauss’s Law
Gauss’s law is a law relating the distribution of electric charge to the resulting electric field.
learning objectives
Describe relationship between the Gauss’s law and the Coulomb’s law
Gauss’s law, also known as Gauss’s flux theorem, is a law relating the distribution of electric charge to the resulting electric field.
17.5.1 https://phys.libretexts.org/@go/page/14547
The law was formulated by Carl Friedrich Gauss (see ) in 1835, but was not published until 1867. It is one of the four Maxwell’s
equations which form the basis of classical electrodynamics, the other three being Gauss’s law for magnetism, Faraday’s law of
induction, and Ampère’s law with Maxwell’s correction.
Carl Friedrich Gauss: Carl Friedrich Gauss (1777–1855), painted by Christian Albrecht Jensen
Gauss’s law can be used to derive Coulomb’s law, and vice versa. Note that since Coulomb’s law only applies to stationary charges,
there is no reason to expect Gauss’s law to hold for moving charges based on this derivation alone. In fact, Gauss’s law does hold
for moving charges, and in this respect Gauss’s law is more general than Coulomb’s law.
In words, Gauss’s law states that: The net outward normal electric flux through any closed surface is proportional to the total
electric charge enclosed within that closed surface.
The law can be expressed mathematically using vector calculus in integral form and differential form, both are equivalent since
they are related by the divergence theorem, also called Gauss’s theorem. Each of these forms in turn can also be expressed two
ways: In terms of a relation between the electric field E and the total electric charge, or in terms of the electric displacement field D
and the free electric charge.
Gauss’s law has a close mathematical similarity with a number of laws in other areas of physics, such as Gauss’s law for
magnetism and Gauss’s law for gravity. In fact, any “inverse-square law” can be formulated in a way similar to Gauss’s law: For
example, Gauss’s law itself is essentially equivalent to the inverse-square Coulomb’s law, and Gauss’s law for gravity is essentially
equivalent to the inverse-square Newton’s law of gravity.
Key Points
If the electric field is uniform, the electric flux passing through a surface of vector area S is Φ = E ⋅ S = ES cos θ .
E
Key Terms
electric field: A region of space around a charged particle, or between two voltages; it exerts a force on charged objects in its
vicinity.
electric displacement field: A vector field that appears in Maxwell’s equations.
electric charge: A quantum number that determines the electromagnetic interactions of some subatomic particles; by
convention, the electron has an electric charge of -1 and the proton +1, and quarks have fractional charge.
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17.6: Applications of Electrostatics
learning objectives
Describe effect of the DNA replication process on the hydrogen bonds
All living organisms contain either DNA or RNA. These molecules contain genetic information that defines (or at least influences)
all of the organism’s characteristics. In both DNA and RNA, electrostatics plays a major role. For the purpose of this atom we
focus specifically on DNA.
DNA is a large macromolecule that (in three-dimensional space) forms the shape of a double helix, as shown in. The “backbone” of
each helix is formed from alternating deoxyribose and phosphate subunits as illustrated in. Each deoxyribose is connected to a base
that extends between the two backbones. Four bases can be found in DNA: adenine, guanine, cytosine and thymine. Their ordering
in sequence determines the genetic characteristics of the parent organism.
Phosphate in DNA: The complex molecules that make up our DNA are held together by a phosphate-deoxyribose backbone, as
shown.
The bases of each strand of DNA are positioned close to one another but not as a result of formal chemical bonds. Electrostatic
interactions, in this case otherwise known as hydrogen bonds, are what hold these bonds together.
In, (in the instance of guanine and cytosine) there are three instances in which hydrogen atoms are attracted to nearby nitrogen and
oxygen atoms (denoted by dashed lines). For adenine and thymine, there are two such hydrogen bonds. These interactions can be
modeled as electrostatic interactions—electrons from the oxygen and nitrogen atoms are negative, and are attracted to the exposed
nuclei (which are positive) of the nearby hydrogen atoms. These hydrogen bonds are relatively weak, but their sum is strong
enough to hold the two DNA strands together.
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Hydrogen bonding between Guanine and Cytosine: In the instance of guanine and cytosine, there are three instances in which
hydrogen atoms are attracted to nearby nitrogen and oxygen atoms (denoted by dashed lines).
When DNA replicates, as it does in cellular reproduction, it is first “unzipped” by any of a number of enzymes known as DNA
Helicases. These enzymes break the electrostatic hydrogen bonds between the two strands.
learning objectives
Describe functions of a photocopier and distinguish the three steps of xerography
A photocopier is a machine that prints copies of documents and images on paper. Currently, most photocopiers use a process called
xerography, which is dry and uses heat in printing.
Different models of photocopiers vary somewhat in their methods, but the first instance of xerography involved use of a
photosensor with a copy camera and a separate processing unit. To this day, xerography is conducted in a five-step process, in
which electrostatics play an important role. This process is diagrammed in with explanations of each step outlined in the text below.
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How a photocopier works: This image describes how a photocopier works.
Charging
In the first step of xerography, a high-voltage device (either a corona wire or charge roller) charges a cylindrical drum. This occurs
by corona discharge, with output limited by a control grid or screen: a negative charge on the wire ionizes the space between the
wire and conductor, so electrons are repelled and pushed onto the conductor. The conductor is put on top of a conducting surface
and is kept at ground potential.
Polarity of charge can be altered depending on desired effects. A positive process is used to produce black on white copies, while a
negative process is used to produce black on white from negative originals.
Exposure
A lamp illuminates the original document. The white areas of the original reflect light onto the surface of the drum, which exhibits
photoconductivity (it becomes a conductor in the presence of light). The areas of the drum exposed to light then discharge to the
ground; the other parts of the drum (which are not conductive, having not been exposed to light) remain negatively charged. Thus,
the result is an electric image on the surface of the drum.
Developing
The toner, a powder used to form images on the paper, is positively charged. When it is placed on the drum, it is attracted to the
negative (black) areas. This develops the image.
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Transfer
The image with toner on the drum is transferred to a piece of paper with more negative charge than the drum.
Fusing
The toner is melted and bonded to the paper by heat and pressure rollers.
It should be noted that the explanation above is typical for modern copiers, but some older copiers use a positive drum and paper,
and a negatively charged toner.
learning objectives
Identify and describe functions of the major parts of a Van de Graaff generator
A Van de Graaff generator is a device used to separate electric charges. Using a moving belt, it can create extremely high potential
differences. In some cases these can reach a value in excess of five megavolts.
Fundamentally, a Van de Graaff generator consists of a flexible dielectric belt (silk is commonly used) running over two metal
pulleys. One of these pulleys is surrounded by a hollow metal sphere. Two electrodes are positioned near the bottom of the lower
pulley and inside the sphere, over the upper pulley. One comb is connected to the sphere, and another is connected to the ground. In
this figure, a high, positive DC potential is applied to the upper roller.
Schematic of a Van de Graaff Generator: Numbers in the diagram indicate: 1) hollow metal sphere; 2) upper electrode; 3) upper
roller (for example an acrylic glass); 4) side of the belt with positive charges; 5) opposite side of the belt with negative charges; 6)
lower roller (metal); 7) lower electrode (ground); 8) spherical device with negative charges, used to discharge the main sphere; 9)
spark produced by the difference of potentials
The separation of charge in a Van de Graaff generator is a complex, multistep process. Due to the electric field around the lower
pulley, the belt receives negative charge as it passes the lower comb. When the belt comes into contact with the upper roller, it
leaves the roller with a negative charge as electrons are transferred. Electrons seep from the belt to the upper comb and then the
terminal, leaving the belt positively charged and the terminal negatively charged. The sphere acts as a Faraday shield, shielding the
upper roller and comb from the electric field produced by charges on the outside of the sphere. Eventually, discharge occurs, and
the belt changes polarity. As the belt continues, a charging current travels via the belt, and the sphere takes on more and more
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negative charge until the rates at which it gains and loses charge (which occurs by leakage and corona discharge) are equal. Final
potential is proportional to the size of the sphere and its distance from the ground.
Key Points
In DNA, hydrogen bonds between base pairs hold two opposite strands together.
The hydrogen bonds that hold DNA together are electrostatic interactions between the nucleus of a hydrogen atom and the
electrons of an oxygen or a nitrogen atom.
When DNA replicates, a helicase enzyme “unzips” the double helix, breaking the hydrogen bonds that hold it together in the
center.
In the first step of xerography, a high-voltage device (either a corona wire or charge roller) charges a cylindrical drum.
In the second step, a lamp illuminates the original document. The white areas of the original reflect light onto the surface of the
drum, which is photoconductive. The areas of the drum exposed to light then discharge to the ground.
Next, the toner, a positively-charged powder used to form images on the paper, is placed on the drum. It is attracted to the
negative (black) areas. This develops the image.
The image with toner on the drum is then transferred to a piece of paper with more negative charge than the drum.
Finally, the toner is melted and bonded to the paper by heat and pressure rollers.
Fundamentally, a Van de Graaff generator consists of a flexible dielectric belt (silk is commonly used) running over two metal
pulleys.
The belt is used to transport charge, isolating positive charge from negative charge.
The maximum potential difference a Van de Graaff generator can create is achieved when the rates at which the sphere (which
holds charge) gains and loses charge are equal. Charge can be lost by leakage and corona discharge.
Key Terms
base: A nucleotide’s nucleobase in the context of a DNA or RNA biopolymer.
photoconductivity: An increase in the electrical conductivity of a material as a result of incident electromagnetic radiation
xerography: a photocopying process in which a negative image formed on an electrically charged plate is transferred as a
positive to paper and thermally fixed
Faraday shield: A Faraday cage or Faraday shield is an enclosure formed by conducting material or by a mesh of such
material. Such an enclosure blocks external static and non-static electric fields.
polarity: the separation, alignment or orientation of something into two opposed poles
corona discharge: an electrical discharge brought on by the ionization of a fluid surrounding a conductor that is electrically
energized
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CHAPTER OVERVIEW
This page titled 18: Electric Potential and Electric Field is shared under a not declared license and was authored, remixed, and/or curated by
Boundless.
1
18.1: Overview
learning objectives
Explain the relationship between the electric potential and the electric field
The relationship between electric potential and field is similar to that between gravitational potential and field in that the potential
is a property of the field describing the action of the field upon an object (see ).
Electric field and potential in one dimension: The presence of an electric field around the static point charge (large red dot)
creates a potential difference, causing the test charge (small red dot) to experience a force and move.
The electric field is like any other vector field—it exerts a force based on a stimulus, and has units of force times inverse stimulus.
In the case of an electric field the stimulus is charge, and thus the units are NC-1. In other words, the electric field is a measure of
force per unit charge.
The electric potential at a point is the quotient of the potential energy of any charged particle at that location divided by the charge
of that particle. Its units are JC-1. Thus, the electric potential is a measure of energy per unit charge.
In terms of units, electric potential and charge are closely related. They share a common factor of inverse Coulombs (C-1), while
force and energy only differ by a factor of distance (energy is the product of force times distance).
Thus, for a uniform field, the relationship between electric field (E), potential difference between points A and B (Δ), and distance
between points A and B (d) is:
Δϕ
E =− (18.1.1)
d
The -1 coefficient arises from repulsion of positive charges: a positive charge will be pushed away from the positively charged
plate, and towards a location of higher-voltage.
The above equation is an algebraic relationship for a uniform field. In a more pure sense, without assuming field uniformity,
electric field is the gradient of the electric potential in the direction of x:
dV
Ex = − (18.1.2)
dx
This can be derived from basic principles. Given that ∆P=W (change in the energy of a charge equals work done on that charge), an
application of the law of conservation of energy, we can replace ∆P and W with other terms. ∆P can be substituted for its definition
as the product of charge (q) and the differential of potential (dV). We can then replace W with its definition as the product of q,
electric field (E), and the differential of distance in the x direction (dx):
qdV = −qExdx (18.1.3)
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Electric Potential Energy and Potential Difference
Electric potential energy results from forces between charges; potential difference is the energy needed to move a charge from point
A to B.
learning objectives
Calculate the potential energy between the charges
Electric potential energy is a type of potential energy that results from Coulomb forces. It is measured in joules and depends on the
positioning of charged particles relative to one another, as well as the magnitude of their respective charges.
The potential energy (UE) between charges q and Q can be calculated as a function of distance between the charges (r):
qQ
UE (r) = (18.1.4)
4π ϵ0 r
If there are three or more charges, the above formula can be modified so that the potential energies between all charges are
summed. Consider, for example, the case involving charges Q1, Q2 and Q3:
1 Q1 Q2 Q2 Q3 Q1 Q3
UE = ( + + ) (18.1.5)
4πϵ0r r12 r23 r13
In this example, r12 represents the distance between Q1 and Q2, r23 represents the distance between Q2 and Q3, and r13represents the
distance between Q1 and Q3. The above formula can be modified for any number of charges.
Potential Difference
Potential difference, or voltage, is the difference in electric potential energy between two points. It is denoted by ∆V and has units
of volts, or joules per Coulomb.
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Electric Potential Difference: A brief overview of electric potential difference and electric potential energy for beginning physics
students.
Voltage denotes the work per unit charge that must be done against a static electric field to move a charge from one point to
another. It may represent a source of energy, or lost, stored or used energy. Voltage also is defined such that negative charges are
pulled towards higher voltages, while positive charges move towards lower voltages. Thus, current in wires flows from higher to
lower voltages.
Potential difference is independent of path taken from one point to another, and may be measured by any of a number of
instruments. These include the voltmeter, the potentiometer, and the oscilloscope. It is most typically measured in circuits, and in
such situations can be calculated using Ohm ‘s Law, which will be covered in a later atom.
Potential difference in a static field: When a charge q moves from point A to point B, the potential difference is independent of
path taken.
learning objectives
Calculate the electric potential created by a charge distribution of constant value
Any charge will create a vector field around itself (known as an electric field ). Electric field is the gradient of potential, which
depends inversely upon distance of a given point of interest from a charge. Placing a second charge in the system (a “test charge”)
results in the two charges experiencing a force (the field’s units are Newtons, a measure of force per Coulomb), causing the charges
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to move relative to one another. It is easiest to model interactions between two charges such that one is considered stationary while
the test charge moves.
As the test charge moves, the potential between it and another charge changes, as does the electric field. The relationship between
potential and field (E) is a differential: electric field is the gradient of potential (V) in the x direction. This can be represented as:
dV
Ex = − (18.1.6)
dx
Equipotential Lines: An isolated point charge Q with its electric field lines (blue) and equipotential lines (green)
Thus, as the test charge is moved in the x direction, the rate of the its change in potential is the value of the electric field.
The instant before the test charge moves, its potential energy is at a maximum, and its kinetic energy is 0. For any charge of
constant value (Q), the potential at a certain distance from it (r) can be calculated from the following equation:
1 Q
VE = (18.1.7)
4πϵ0 r
where ε0 is the electric constant, otherwise known as permittivity of free space. Moving towards and away from the charge results
in change of potential; the relationship between distance and potential is inverse.
For one point charge, potential will be constant for all points a certain radial distance away. Multiple points of the same potential
are known as equipotential. In the case of fields created by a single point charge, all points on any circle centered around the point
charge will be equipotential, as illustrated in.
shows that when multiple charges create a field, the equipotential lines become irregularly shaped. This is due to the fact that the
fields created by each charge overlap, thus potential is increased at any point relative to that which would have arisen from one or
the other charge.
learning objectives
Determine the electric potential within and outside a charged conductor
When a conductor becomes charged, that charge distributes across its surface until electrostatic equilibrium is reached. Its surface
is equipotential.
All points within a charged conductor experience an electric field of 0. This is because field lines from charges on the surface of the
conductor oppose one another equally. However, having the electric field equal to zero at all points within a conductor, the electric
potential within a conductor is not necessarily equal to zero for all points within that same conductor. This can be proven by
relating electric field and potential.
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Electrical Charge at a Sharp Point of a Conductor: Repulsive forces towards the more sharply curved surface on the right aim
more outward than along the surface of the conductor.
Given that work is the difference in final and initial potential energies (∆U), we can relate this difference to the dot product of force
at every infinitesimal distance l along the path between the points within the conductor:
\Delta \mathrm { U } = - \int _ { \mathrm { i } } ^ { f } \vec { \mathrm { F } } \cdot \mathrm { d } \vec { }
This is the equation for work, with ∆U substituted in place of W. Rewriting U as the product of charge (q) and potential difference
(V), and force as the product of charge and electric field (E), we can assert:
f
⃗
Δ(qV) = − ∫ (qE) ⋅ dl (18.1.8)
i
Dividing both sides by the common term of q, we simplify the equation to:
f
⃗ ⃗
ΔV = − ∫ E ⋅ dl (18.1.9)
i
Thus we can conclude that, given that the electric field is constantly 0 for any location within the charged conductor, the potential
difference in that same volume needs to be constant and equal to 0.
On the other hand, for points outside a conductor, potential is nonzero and can be defined by the very same equation, according to
field and distance from the conductor.
learning objectives
Describe properties and approximations of the uniform electric field
A uniform field is that in which the electric field is constant throughout. Just like the so-called “frictionless surface” in mechanics,
the uniform field is an ideal but unreal situation that makes for simpler calculations. Equations involving non-uniform electric
fields require use of differential calculus.
Uniformity in an electric field can be approximated by placing two conducting plates parallel to one another and creating a
potential difference between them. In such a case there will be slight variations in the field near its edges, but it will be
approximately constant throughout every other area.
The equation for magnitude of a uniform electric field is:
−Δϕ
E = (18.1.11)
d
where E is the field, Δ is the potential difference between the plates, and d is the distance between the plates. The coefficient of -1
arises from the fact that positive charges repel, and thus a positive charge will be pushed away from the positive plate and in a
direction opposite that of the increasing voltage.
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Uniformity of an electric field allows for simple calculation of work performed when a test charge is moved across it. For the case
of a positive charge q to be moved from a point A with a certain potential (V1) to a point B with another potential (V2), that
equation is:
W = −q (V2 − V1 ) (18.1.12)
The difference (V2-V1) can also be represented as ∆V or VAB. In uniform fields it is also simple to relate ∆V to field strength and
distance (d) between points A and B:
VAB = Ed (18.1.13)
Relationships within a uniform electric field: In this image, Work (W), field strength (E), and potential difference (∆V) are
defined for points A and B within the constructs of a uniform potential field between the positive and negative plates.
Energy Conservation
Energy is conserved in the movement of a charged particle through an electric field, as it is in every other physical situation.
learning objectives
Formulate energy conservation principle for a charged particle in an electric field
Energy is conserved in the movement of a charged particle through an electric field, as it is in every other physical situation. This
phenomenon can be expressed as the equality of summed kinetic (Ekin) and electric potential (Eel) energies:
(Ekin + Eel ) initial = (Ekin + Eel )f inal (18.1.14)
Given a stationary test charge in a certain location, an applied electric field will cause the charge to move to one end or the other,
depending on the charge (positive test charges will move in the direction of the field; negative charges will move in the opposite
direction). In all cases, a charge will naturally move from an area of higher potential energy to an area of lower potential energy.
At the instant at which the field is applied, the motionless test charge has 0 kinetic energy, and its electric potential energy is at a
maximum. After that moment, the charge accelerates, and its kinetic energy (from motion) increases as its potential energy
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decreases. Throughout this time, the sum of potential and kinetic energies remains constant.
Another way to express the previous equation is:
1 2
1 2
( mv + U) initial = ( mv + U)f inal (18.1.15)
2 2
where m and v are the mass and velocity of the electron, respectively, and U is the electric potential energy. U can be calculated as
follows:
q0 q
U = q0 V = k (18.1.16)
r
where V is the potential difference, k is a constant, q0 is a test charge, q is another charge, and r is the distance between the charges.
The terms involved in the formula for conservation of energy can be rewritten in many ways, but all expressions are based on the
simple premise of equating the initial and final sums of kinetic and potential energy.
Similarities Between Activity of Gravitational and Electric Fields on an Object: The charge, +q, is moved down the electric
field in the same way that the object, m, is moved down the hill. In both instances, the particle in motion goes from a higher to a
lower potential energy state.
The Electron-Volt
The electron volt is a unit of energy useful in the physics of elementary charges and electricity.
learning objectives
Convert between electron volts and SI units of energy
Overview
The electron volt, symbolized as eV and sometimes written as electronvolt, is a unit of energy useful in the physics of elementary
charges and electricity.
The electron volt is defined as the amount of energy gained or lost by the charge of an electron moved across a one-volt electric
potential difference. As such, it is equal to the product of one volt (1 J/C) and one elementary charge, giving it a value in joules
approximately equal to 1.602×10-19 J.
Not an SI unit in itself, the electron volt became useful through experimentation. Scientists working with electrostatic particle
accelerators commonly used the relationship between energy (E), charge (q), and potential difference (V) in their work:
E = qV (18.1.17)
All calculations of energy from the above equation were quantized as multiples of the elementary charge, q, for a given voltage,
and thus arose the common usage of the electron volt as a unit of measurement.
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Momentum
Both electron volts and momentum are measures of energy, and the two are related in high-energy physics. Applying a potential
difference to an electron gives it energy, which manifests itself in motion of the electron through it. Given that the electron has both
mass and velocity, it has momentum. Dividing electron volts by a constant with units of velocity results in a momentum.
Mass
Given that mass is equivalent to energy, the electron volt can measure mass. In particle physics, the equation E=mc2 can be
rearranged to solve for mass:
E
m = (18.1.18)
c2
Wavelength
The energy E, frequency v, and wavelength λ of a photon are related by
hc
E(eV) = hv = (18.1.19)
λ
where h is the Planck constant and c is the speed of light. Thus, a photon with a wavelength of 532 nm (green light) would have an
energy of approximately 2.33 eV. Similarly, 1 eV would correspond to an infrared photon of wavelength 1240 nm, and so on.
Energy of Photons in the Visible Spectrum: Relationship between wavelength and energy expressed in electron volts.
Temperature
In plasma physics, the electron volt can be used as a unit of temperature. To convert to Kelvins, simply divide the value of 1 eV (in
Joules) by the Boltzmann constant (1.3806505(24)×10-23 J/K).
Dipole Moments
The electric dipole moment is a measure of polarity in a system.
learning objectives
Relate the electric dipole moment to the polarity in a system
The electric dipole moment is a measure of polarity, which is the separation of positive and negative charges in a system. It is
measured in units of Coulomb-meters (C m). There are many different types of dipole moments, including electric dipole moments,
magnetic dipole moments, and topological dipole moments.
Among the subset of electric dipole moments are transition dipole moments, molecular dipole moments, bond dipole moments, and
electron electric dipole moments. For the purposes of this atom we will focus on a broad overview of electric dipole moment in
static situations.
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Molecular Dipole Moment in Water: This water (H2O) molecule has a high density of electrons (denoted by the red shading) near
the red O atom. Closer to the white H atoms, there is a low density of electrons. Therefore, the molecule is is a dipole, with
negativity near the O and positivity closer to the H atoms.
Definition
Fundamentally, for the case of point charges with values +q and -q, an electric dipole moment (p) can be defined as the vector
product of the charges and the displacement vector d:
p = qd (18.1.20)
The displacement vector is the vector with a magnitude equal to the distance between the charges and a direction pointing from the
negative charge to the positive charge. It is essentially interchangeable with the “radius” variable in many other equations (such as
those determining gravitational and electrostatic forces ), except that it includes the factor of direction.
Torque
All dipoles will experience a torsional force, or torque, when they are placed in external electric fields. This torque rotates the
dipole to align it with the field. It is brought on by the need to minimize potential energy. Torque (τ) can be calculated as the cross
product of the electric dipole moment and the electric field (E), assuming that E is spatially uniform:
τ = p×E (18.1.21)
Key Points
The electric field is a measure of force per unit charge; the electric potential is a measure of energy per unit charge.
For a uniform field, the relationship between electric field (E), potential difference between points A and B (Δ), and distance
Δϕ
between points A and B (d) is: E = − If the field is not uniform, calculus is required to solve.
d
Potential is a property of the field that describes the action of the field upon an object.
Electric potential energy is a type of potential energy that results from Coulomb forces. The potential energy (UE) between
qQ
charges q and Q can be calculated as a function of distance between the charges (r): U (r) = . E
4πϵ0 r
The formula for potential energy can be modified for potential between many charges, so long as the interactions of each charge
with every other charge in the system are considered. For example, potential between three charges can be solved using the
Q1 Q2 Q2 Q3 Q1 Q3
following formula: U E =
1
4πϵ0 r
(
r12
+
r23
+
r13
) .
Potential difference, or voltage, is the difference in electric potential energy between two points. It is denoted by ∆V and has
units of volts, or joules per Coulomb.
For any charge of constant value (Q), the potential (VE) at a certain distance from it (r) can be calculated from the equation:
Q
VE =
1
4πϵ0 r
Where ε is the electric constant, otherwise known as permittivity of free space.
0
For one point charge, potential will be constant for all points a certain radial distance away. Multiple points of the same
potential are known as equipotential.
When multiple charges create a field, the equipotential lines become irregularly shaped. This is because the fields created by
each charge overlap, thus the potential is increased at any point relative to that which would have arisen from one or the other
charge.
f
Electric potential (∆V) and field (E) are related according to the integral: ΔV = − ∫ E⃗ ⋅ dl ⃗ where l is the distance between
i
18.1.9 https://phys.libretexts.org/@go/page/14550
In uniform fields it is also simple to calculate potential difference: V = Ed In this case, field strength is E, and distance
AB
Key Terms
electric field: A region of space around a charged particle, or between two voltages; it exerts a force on charged objects in its
vicinity.
electric potential: The potential energy per unit charge at a point in a static electric field; voltage.
coulomb: In the International System of Units, the derived unit of electric charge; the amount of electric charge carried by a
current of 1 ampere flowing for 1 second. Symbol: C
potential energy: The energy an object has because of its position (in a gravitational or electric field) or its condition (as a
stretched or compressed spring, as a chemical reactant, or by having rest mass)
equipotential: A region whose every point has the same potential.
radial: Moving along a radius.
electric field: A region of space around a charged particle, or between two voltages; it exerts a force on charged objects in its
vicinity.
work: A measure of energy expended in moving an object; most commonly, force times displacement. No work is done if the
object does not move.
potential difference: The difference in potential energy between two points in an electric field; the difference in charge
between two points in an electrical circuit; voltage.
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
particle accelerator: A device that accelerates electrically charged particles to extremely high speeds, for the purpose of
inducing high-energy reactions or producing high-energy radiation.
electron volt: A unit for measuring the energy of subatomic particles; the energy equal to that attained by an electron moving
through a potential difference of one volt. Equivalent to 1.6022 x 10-19 joules.
dipole moment: The vector product of the charge on either pole of a dipole and the distance separating them.
vector: A directed quantity, one with both magnitude and direction; the between two points.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
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18.2: Equipotential Surfaces and Lines
learning objectives
Compare conductivity and resistivity of an ideal conductor, commenting on the presence of ideal conductors in nature
An ideal conductor is one that exists only in the world of theory. It is one that has certain “ideal” properties that make calculations
of electric potential and field, and other properties simple to perform; it essentially eliminates the need for “corrections” to account
for slight deviations from the so-called “ideal. ” These ideal properties are:
σ =∞ (18.2.1)
Conductivity (σ) is the inverse of resistance, measured in units of current per potential difference. In an ideal conductor, the
material’s conductivity is infinite, and its resistance approaches 0. This means that, ideally, it needs a minute amount of voltage
(potential difference) to carry an extremely high amperage (current). The principle of near-zero resistance is akin to that of
frictionless surfaces: t. Theoretically, with the slightest force (voltage), an object (current) on a frictionless surface (zero-resistance
conductor) can proceed without restriction.
Etan = Dtan = 0 (18.2.2)
The electric field (Etan) and electric flux density (Dtan) tangential to the surface of a conductor must be equal to 0. This is because
any such field or flux that is tangential to the surface of the conductor must also exist inside the conductor, which by definition
touches the tangential field or density at one point.
If an electric field exists inside the metal, there must be a drop in voltage between any two points along the surface of the metal. In
a perfect conductor this drop should not exist because it implies a less-than-infinite conductivity.
Electric flux density normal to the conductor’s surface is equal to surface charge density. Essentially this means that the
conductor’s charge exists on its surface, not in its interior.
→
E int =0 (18.2.4)
This means that the electric field inside a perfect conductor is 0. The charges along the surface all act equally and opposite to one
another, and their sum at any point is equal to 0.
Φsurf ace = constant (18.2.5)
Charge distribution may vary depending on shape, but the potential over the surface of an ideal conductor is, at electrostatic
equilibrium, constant throughout.
18.2.1 https://phys.libretexts.org/@go/page/14551
Charge Distribution on a Conductor with an Irregular Surface: Curvature causes electric field lines to extend such that they
further distance themselves from one another with increasing distance from the conductor surface. As such, charges (and field
lines) aggregate around areas of curvature.
learning objectives
Give examples of the electric potentials in human body
Electric potentials are not limited in function to inorganic processes. In fact, they can be commonly seen in living organisms. In
humans, they are seen in cell membranes and nerve impulses in particular.
Cell Membranes
Cell membranes are only semipermeable; water can freely travel in and out, but ions can be selectively admitted passage across
them. As a result, a cell can contain a concentration of a given ion that differs from that which exists outside. Thus, a potential,
called the resting potential, is created on either side of the membrane.
Typical ions used to generate resting potential include potassium, chloride, and bicarbonate.
Resting membrane potential is approximately -95 mV in skeletal muscle cells, -60 mV in smooth muscle cells, -80 to -90 mV in
astroglia, and -60 to -70 mV in neurons.
Potentials can change as ions move across the cell membrane. This can occur passively, as ions diffuse through ion channels in the
membrane. No energy is required for this to occur, and therefore ions can only move from areas of higher concentration to those of
lower concentration.
18.2.2 https://phys.libretexts.org/@go/page/14551
Active transport of ions across a cell membrane is also a possibility. This involves ion pumps using energy to push an ion from an
area of lower concentration to one of higher concentration.
Nerve Impulses
When the brain decides on an action, it sends an impulse that cascades to the extremity where a muscle contracts.
Neurons receive an impulse at the dendrites. This impulse is passed through the axon, a long extension of the cell, in the form of an
electrical potential created by differing concentrations of sodium and potassium ions on either side of a membrane in the axon.
The Neuron: Neurons receive an impulse at the dendrites. This impulse is passed through the axon, a long extension of the cell, in
the form of an electrical potential created by differing concentrations of sodium and potassium ions on either side of a membrane in
the axon.
When the signal reaches the end of the axon, neurotransmitters are released, which then are received by the dendrites of the next
neuron. The next neuron repeats the process outlined above.
Equipotential Lines
Equipotential lines depict one-dimensional regions in which the electric potential created by one or more nearby charges is
constant.
learning objectives
Describe the shape of the equipotential lines for several charge configurations
Equipotential lines depict one-dimensional regions in which the electric potential created by one or more nearby charges has a
constant value. This means that if a charge is at any point on a given equipotential line, no work will be required to move it from
one point to another on that same line.
18.2.3 https://phys.libretexts.org/@go/page/14551
High School Physics - Parallel Plates and Equipotent…
Equipotent…
Parallel Plates and Equipotential Lines: A brief overview of parallel plates and equipotential lines from the viewpoint of
electrostatics.
Equipotential lines may be straight, curved, or irregularly shaped, depending on the orientation of charges that give rise to them.
Since they are located radially around a charged body, they are perpendicular to electric field lines, which extend radially from the
center of a charged body.
The radial dependence means that at any point a certain distance from the point charge, potential will be the same. Therefore,
equipotential lines for a single point charge are circular, with the point charge at the center.
18.2.4 https://phys.libretexts.org/@go/page/14551
Equipotential Lines: An isolated point charge Q with its electric field lines (blue) and equipotential lines (green)
Equipotential lines with multiple charges: Equipotential lines become irregularly shaped when there are multiple charges vicinal
to one another. At a point between the charges, a test charge may “feel” the effects of both charges.
Continuous Charge
If charges are distributed across two conductor plates in static equilibrium, in which charges are continuous and distributed in a
straight line, the equipotential lines will be approximately straight. This is because the continuity of charges results in a continuity
of action upon any point a certain distance from either plate.
18.2.5 https://phys.libretexts.org/@go/page/14551
Equipotential Lines Between Two Plates: When charges are lined up and continuous on conducting plates, equipotential lines are
straight between them. The only exception is a curving of the lines near the edges of the conductor plates.
This continuity is broken towards the ends of the plates, however, which causes curvature in these areas. This curvature is known
as “edge effects. ”
Key Points
In an ideal conductor, conductivity is assumed to be infinite (as resistivity approaches 0).
The electric field (Etan) and electric flux density (Dtan) tangential to the surface of a conductor must be equal to 0.
Electric flux density normal to the conductor’s surface is equal to surface charge density. Essentially this means that the
conductor’s charge exists on its surface, not in its interior.
Electric field inside a perfect conductor is 0. The charges along the surface all act equally and opposite to one another, and their
sum at any point is equal to 0.
Charge distribution may vary depending on shape, but the potential over the surface of an ideal conductor is, at electrostatic
equilibrium, constant throughout.
The resting potential is a potential created by an imbalance of ions on either side of a cell membrane.
Resting potential can be altered by passive diffusion (which requires no energy ) or active transport of ions across a cell
membrane.
When the brain decides on an action, it sends an impulse that cascades to the extremity where a muscle contracts. Along the
axon of a neuron, this impulse manifests itself in a potential created by an imbalance of sodium and potassium ions across a
membrane.
For a single, isolated point charge, potential is inversely dependent upon radial distance from the charge. Therefore,
equipotential lines for a single point charge are circular, with the point charge at the center.
When multiple, discrete charges interact, their fields overlap, meaning their potentials combine. This results in the skewing of
equipotential lines in areas that are close enough to both charges to “feel” the effects of both fields.
If charges are distributed across two conductor plates in static equilibrium, in which charges are continuous and distributed in a
straight line, the equipotential lines will be approximately straight.
18.2.6 https://phys.libretexts.org/@go/page/14551
Key Terms
conductor: A material which contains movable electric charges.
flux density: A measure of rate of flow of a fluid, particles or energy per unit area.
neuron: A cell of the nervous system, which conducts nerve impulses; consisting of an axon and several dendrites. Neurons are
connected by synapses.
axon: A nerve fiber which is a long, slender projection of a nerve cell, and which conducts nerve impulses away from the body
of the cell to a synapse.
cell membrane: The semipermeable membrane that surrounds the cytoplasm of a cell.
equipotential: A region whose every point has the same potential.
static equilibrium: the physical state in which all components of a system are at rest and the net force is equal to zero
throughout the system
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18.3: Point Charge
learning objectives
Express the electric potential generated by a single point charge in a form of equation
The electric potential tells you how much potential energy a single point charge at a given location will have. The electric potential
at a point is equal to the electric potential energy (measured in joules) of any charged particle at that location divided by the charge
(measured in coulombs) of the particle. Since the charge of the test particle has been divided out, the electric potential is a
“property” related only to the electric field itself and not the test particle. Another way of saying this is that because PE is
dependent on q, the q in the above equation will cancel out, so V is not dependent on q.
The potential difference between two points ΔV is often called the voltage and is given by
ΔPE
ΔV = VB − VA = (18.3.2)
q
Point Charges
Point charges, such as electrons, are among the fundamental building blocks of matter. Furthermore, spherical charge distributions
(like on a metal sphere, see figure below) create external electric fields exactly like a point charge. The electric potential due to a
point charge is, thus, a case we need to consider. Using calculus to find the work needed to move a test charge q from a large
distance away to a distance of r from a point charge Q, and noting the connection between work and potential (W=–qΔV), it can be
shown that the electric potential V of a point charge is
kQ
V =
r
(point charge)
where k is a constant equal to 9.0×109 N⋅m2/C2 .
18.3.1 https://phys.libretexts.org/@go/page/14552
Van de Graaff Generator: The voltage of this demonstration Van de Graaff generator is measured between the charged sphere and
ground. Earth’s potential is taken to be zero as a reference. The potential of the charged conducting sphere is the same as that of an
equal point charge at its center.
The potential at infinity is chosen to be zero. Thus V for a point charge decreases with distance, whereas E for a point charge
decreases with distance squared:
F kQ
E = = (18.3.3)
2
q r
The electric potential is a scalar while the electric field is a vector. Note the symmetry between electric potential and gravitational
potential – both drop off as a function of distance to the first power, while both the electric and gravitational fields drop off as a
function of distance to the second power.
learning objectives
Explain how the total electric potential due to a system of point charges is found
We’ve also seen that the electric potential due to a point charge is
where k is a constant equal to 9.0×109 N⋅m2/C2. The equation for the electric potential of a point charge looks similar to the
equation for the electric field generated for a point particle
18.3.2 https://phys.libretexts.org/@go/page/14552
F kQ
E = = (18.3.5)
2
q r
with the difference that the electric field drops off with the square of the distance while the potential drops off linearly with
distance. This is analogous to the relationship between the gravitational field and the gravitational potential.
Superposition of Electric Potential: The electric potential at point L is the sum of voltages from each point charge (scalars).
Recall that the electric potential V is a scalar and has no direction, whereas the electric field E is a vector. To find the voltage due to
a combination of point charges, you add the individual voltages as numbers. So for example, in the figure above the electric
potential at point L is the sum of the potential contributions from charges Q1, Q2, Q3, Q4, and Q5 so that
Q1 Q2 Q3 Q4 Q5
VL = k [ + + + + ] (18.3.6)
d1 d2 d3 d4 d5
To find the total electric field, you must add the individual fields as vectors, taking magnitude and direction into account. This is
consistent with the fact that V is closely associated with energy, a scalar, whereas E is closely associated with force, a vector.
The summing of all voltage contributions to find the total potential field is called the superposition of electric potential. Summing
voltages rather than summing the electric simplifies calculations significantly, since addition of potential scalar fields is much
easier than addition of the electric vector fields. Note that there are cases where you might need to sum potential contributions from
sources other than point charges; however, that is beyond the scope of this section.
Key Points
Recall that the electric potential is defined as the potential energy per unit charge, i.e. V = PE
q
.
The potential difference between two points ΔV is often called the voltage and is given b ΔV = V − V = . The B A
ΔPE
The electric potential V is a scalar and has no direction, whereas the electric field E is a vector.
To find the voltage due to a combination of point charges, you add the individual voltages as numbers. So for example, in the
electric potential at point L is the sum of the potential contributions from charges Q1, Q2, Q3, Q4, and Q5 so that
Q1 Q2 Q3 Q4 Q5
VL = k [
d1
+
d2
+
d3
+
d4
+
d5
] .
To find the total electric field, you must add the individual fields as vectors, taking magnitude and direction into account. This is
consistent with the fact that V is closely associated with energy, a scalar, whereas E is closely associated with force, a vector.
The summing of all voltage contributions to find the total potential field is called the superposition of electric potential. It is
much easier to sum scalars than vectors, so often the preferred method for solving problems with electric fields involves the
summing of voltages.
18.3.3 https://phys.libretexts.org/@go/page/14552
Key Terms
electric potential: The potential energy per unit charge at a point in a static electric field; voltage.
voltage: The amount of electrostatic potential between two points in space.
vector: A directed quantity, one with both magnitude and direction; the between two points.
scalar: A quantity that has magnitude but not direction; compare vector.
superposition: The summing of two or more field contributions occupying the same space.
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18.4: Capacitors and Dielectrics
learning objectives
Express the relationship between the capacitance, charge of an object, and potential difference in the form of equation
Capacitance is the measure of an object’s ability to store electric charge. Any body capable of being charged in any way has a value
of capacitance.
The unit of capacitance is known as the Farad (F), which can be adjusted into subunits (the millifarad (mF), for example) for ease
of working in practical orders of magnitude. The Farad can be equated to many quotients of units, including JV-2, WsV-2, CV-1,
and C2J-1.
The most common capacitor is known as a parallel-plate capacitor which involves two separate conductor plates separated from
one another by a dielectric. Capacitance (C) can be calculated as a function of charge an object can store (q) and potential
difference (V) between the two plates:
Parallel-Plate Capacitor: The dielectric prevents charge flow from one plate to the other.
q
C = (18.4.1)
V
Ultimately, in such a capacitor, q depends on the surface area (A) of the conductor plates, while V depends on the distance (d)
between the plates and the permittivity (εr) of the dielectric between them. For a parallel-plate capacitor, this equation can be used
to calculate capacitance:
A
C = ϵr ϵ0 (18.4.2)
d
Where ε0 is the electric constant. The product of length and height of the plates can be substituted in place of A.
In storing charge, capacitors also store potential energy, which is equal to the work (W) required to charge them. For a capacitor
with plates holding charges of +q and -q, this can be calculated:
Q 2
q CV
Wcharging = ∫ dq = = Wstored (18.4.3)
0
C 2
Thus, either through calculus or algebraically (if C and V are known), stored energy (Wstored) can be calculated. In a parallel-plate
capacitor, this can be simplified to:
2
ϵϵ0 A V
Wstored = (18.4.4)
2d
18.4.1 https://phys.libretexts.org/@go/page/14553
learning objectives
Describe the behavior of the dielectric material in a capacitor’s electric field
In order for a capacitor to hold charge, there must be an interruption of a circuit between its two sides. This interruption can come
in the form of a vacuum (the absence of any matter) or a dielectric (an insulator).
When a dielectric is used, the material between the parallel plates of the capacitor will polarize. The part near the positive end of
the capacitor will have an excess of negative charge, and the part near the negative end of the capacitor will have an excess of
positive charge. This redistribution of charge in the dielectric will thus create an electric field opposing the field created by the
capacitor.
Diagram of a Parallel-Plate Capacitor: Charges in the dielectric material line up to oppose the charges of each plate of the
capacitor. An electric field is created between the plates of the capacitor as charge builds on each plate.
Therefore, the net field created by the capacitor will be partially decreased, as will the potential difference across it, by the
dielectric. On the other hand, the dielectric prevents the plates of the capacitor from coming into direct contact (which would render
the capacitor useless). If it has a high permittivity, it also increases the capacitance for any given voltage. The capacitance for a
parallel-plate capacitor is given by:
c=ϵAdc=ϵAd
where ε is the permittivity, A is the area of the capacitor plates (assuming both are the same size and shape), and d is the thickness
of the dielectric.
Any insulator can be used as a dielectric, but the materials most commonly used are selected for their ability to resist ionization.
The more resistant a material is to ionization, the more tolerance it has for operating at higher voltages. Eventually every material
has a “dielectric breakdown point,” at which the potential difference becomes too high for it to insulate, and it ionizes and permits
the passage of current.
Parallel-Plate Capacitor
The parallel-plate capacitor is one that includes two conductor plates, each connected to wires, separated from one another by a thin
space.
learning objectives
Calculate the maximum storable energy in a parallel-plate capacitor
One of the most commonly used capacitors in industry and in the academic setting is the parallel-plate capacitor. This is a capacitor
that includes two conductor plates, each connected to wires, separated from one another by a thin space. Between them can be a
vacuum or a dielectric material, but not a conductor.
18.4.2 https://phys.libretexts.org/@go/page/14553
Parallel-Plate Capacitor: In a capacitor, the opposite plates take on opposite charges. The dielectric ensures that the charges are
separated and do not transfer from one plate to the other.
The purpose of a capacitor is to store charge, and in a parallel-plate capacitor one plate will take on an excess of positive charge
while the other becomes more negative.
Assuming the plates extend uniformly over an area of A and hold ± Q charge, their charge density is ±, where ρ=Q/A. Assuming
that the dimensions of length and width for the plates are significantly greater than the distance (d) between them, the electric field
(E) near the center of the plates can be calculated by:
ρ
E = (18.4.5)
ϵ
Potential (V) between the plates can be calculated from the line integral of the electric field (E):
d
V =∫ Edz (18.4.6)
0
where z is the axis perpendicular to both plates. Through simplification and substitution, this integral can be changed to:
ρd Qd
V = = (18.4.7)
ϵ ϵA
Accordingly, capacitance is greatest in devices with high permittivity, large plate area, and minimal separation between the plates.
The maximum energy (U) a capacitor can store can be calculated as a function of Ud, the dielectric strength per distance, as well as
capacitor’s voltage (V) at its breakdown limit (the maximum voltage before the dielectric ionizes and no longer operates as an
insulator):
2 2 2
CV ϵA(Udd ) ϵAdU
d
U = = = (18.4.9)
2 2d 2
learning objectives
Calculate the total capacitance for the capacitors connected in series and in parallel
Like any other form of electrical circuitry device, capacitors can be used in combination in circuits. These combinations can be in
series (in which multiple capacitors can be found along the same path of wire) and in parallel (in which multiple capacitors can be
found along different paths of wire).
Capacitors in Series
Like in the case of resistors in parallel, the reciprocal of the circuit’s total capacitance is equal to the sum of the reciprocals of the
capacitance of each individual capacitor:
18.4.3 https://phys.libretexts.org/@go/page/14553
Capacitors in Series: This image depicts capacitors C1, C2 and so on until Cn in a series.
1 1 1 1
= + +… + (18.4.10)
Ctotal C1 C2 Cn
Parallel Capacitors
Total capacitance for a circuit involving several capacitors in parallel (and none in series) can be found by simply summing the
individual capacitances of each individual capacitor.
Parallel Capacitors: This image depicts capacitors C1, C2, and so on until Cn in parallel.
Ctotal = C1 + C2 + … + Cn (18.4.12)
Capacitors in Series and in Parallel: The initial problem can be simplified by finding the capacitance of the series, then using it
as part of the parallel calculation.
The circuit shown in (a) contains C1 and C2 in series. However, these are both in parallel with C3. If we find the capacitance for the
series including C1 and C2, we can treat that total as that from a single capacitor (b). This value can be calculated as approximately
equal to 0.83 μF.
With effectively two capacitors left in parallel, we can add their respective capacitances (c) to find the total capacitance for the
circuit. This sum is approximately 8.83 μF.
18.4.4 https://phys.libretexts.org/@go/page/14553
learning objectives
Identify conditions that can lead to a dielectric breakdown and its effect on materials
Dielectric breakdown (illustrated in ) is the phenomenon in which a dielectric loses its ability to insulate, and instead becomes a
conductor. Dielectrics are commonly used either to isolate conductors from a variable external environment (e.g., as coating for
electrical wires) or to isolate conductors from one another (e.g., between plates of a parallel-plate capacitor). In all applications,
they are selected for their ability to act as insulators. By definition, an insulator is unable to conduct electricity. Under certain
conditions, however, a material that is an insulator can become a conductor.
Eventually, exposing any insulator to increasing voltage will result in the insulator becoming conductive. This point (the minimum
voltage for the insulator to become a conductor) is known as the breakdown voltage. Breakdown is more of a rough concept than
an exact science. A material’s breakdown voltage cannot be precisely defined. As a failure, there is a probabilistic element and thus
a dielectric may experience a breakdown at any of a range of voltages. Additionally, the nature of the voltage used to induce
breakdown must be considered. Short pulses can be used in stress testing to resemble lightning strikes, as could a continuous
applied voltage.
However, for the case of a gas being used as a dielectric, the following equation has been proven to be rather reliable in predicting
breakdown voltage (Vb):
Bpd
Vb = (18.4.13)
1
ln Apd − ln(ln(1 + ))
γs e
where A and B are constants that depend on the surrounding gas, p is the pressure of the surrounding gas, d is distance between the
electrodes (in cm) and γse is the secondary electron emission coefficient. Gaseous dielectrics commonly experience breakdown in
nature (the phenomenon of lightning is the most common example).
Dielectric breakdown of plexiglas: The treelike pattern in the plexiglas stems from the root of the breakdown. Current is dispersed
in many different directions, creating different stems.
Key Points
The unit of capacitance is known as the farad (F), which can be equated to many quotients of units, including JV-2, WsV-2, CV-
1, and C2J-1.
18.4.5 https://phys.libretexts.org/@go/page/14553
Capacitance (C) can be calculated as a function of charge an object can store (q) and potential difference (V) between the two
q
plates: C = . Q depends on the surface area of the conductor plates, while V depends on the distance between the plates and
V
with plates holding charges of +q and -q, this can be calculated: W = . The above can be equated with the work
stored
CV
(assuming both are the same size and shape), and d is the thickness of the dielectric.
Any insulator can be used as a dielectric, but the materials most commonly used are selected for their ability to resist ionization.
The more resistant a material is to ionization, the more tolerance it has for operating at higher voltages.
Assuming the plates extend uniformly over an area of A and hold ± Q charge, their charge density is ±, where ρ=Q/A.
Assuming that the dimensions of length and width for the plates are significantly greater than the distance (d) between them,
ρ
E = can be used to calculate the electric field (E) near the center of the plates. In this equation, ε represents permittivity.
ϵ
ρd Qd
V =
ϵ
= can be used to calculate the potential between the plates.
ϵA
C =
ϵA
d
can be found from the previous equation, adjusting the terms to solve for capacitance (C).
2 2
2
ϵA( Udd ) ϵAdU
U =
CV
2
= solves for the maximum storable energy in a parallel-plate capacitor (U) as a function of Ud, the
2d
=
2
d
dielectric strength per distance as well as capacitor’s voltage (V) at its breakdown limit.
1
=
Ct ot al
+ +… +
C1
1
Capacitors in series follow the law of reciprocals; the reciprocal of the circuit ‘s total capacitance
1
C2
1
Cn
is equal to the sum of the reciprocals of the capacitances of each individual capacitor.
C total= C +C +… +C 1 For capacitors in parallel, summing the capacitances of individual capacitors affords the total
2 n
Key Terms
dielectric: An electrically insulating or nonconducting material considered for its electric susceptibility (i.e., its property of
polarization when exposed to an external electric field).
capacitance: The property of an electric circuit or its element that permits it to store charge, defined as the ratio of stored
charge to potential over that element or circuit (Q/V); SI unit: farad (F).
capacitor: An electronic component capable of storing an electric charge, especially one consisting of two conductors separated
by a dielectric.
permittivity: A property of a dielectric medium that determines the forces that electric charges placed in the medium exert on
each other.
circuit: A pathway of electric current composed of individual electronic components, such as resistors, transistors, capacitors,
inductors and diodes, connected by conductive wires or traces through which electric current can flow. T
conductor: A material which contains movable electric charges.
breakdown: A failure, particularly mechanical; something that has failed.
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18.5: Applications
learning objectives
Identify the primary components of a cathode ray tube and describe the use of cathode ray tubes
Color Cathode Ray Tube: Cutaway rendering of a color CRT: 1) Three Electron guns (for red, green, and blue phosphor dots) 2)
Electron beams 3) Focusing coils 4) Deflection coils 5) Anode connection 6) Mask for separating beams for red, green, and blue
part of displayed image 7) Phosphor layer with red, green, and blue zones 8) Close-up of the phosphor-coated inner side of the
screen
The CRT uses an evacuated glass envelope which is large, deep (i.e., long from front screen face to rear end), fairly heavy, and
relatively fragile. As a matter of safety, the face is typically made of thick lead glass so as to be highly shatter-resistant and to block
most X-ray emissions, particularly if the CRT is used in a consumer product.
Phosphors
The phosphors in a CRT’s screen are the materials that directly produce the photons generated by the CRT. These phosphors are
struck by incoming electrons from the electron gun, absorb energy, and then re-emit some or all of that energy in the form of light
(this process is called phosphorescence). By varying the type of phosphor used, one may vary the wavelength of light emitted by
the phosphor when excited. For example, black and white TV screens use one type of phosphor, while color TVs use three (blue,
red, and green). Early computer terminal monitors used only green phosphors.
CRT Devices
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Monochrome Computer CRT Monitor: Monochrome monitor – this CRT uses only one type of phosphor.
Oscilloscope
An oscilloscope is a device that measures and displays voltages as a time versus voltage graph. The voltage difference between the
positive and negative probe leads is measured, buffered, and displayed on the screen as a continuous curve. Oscilloscopes are
generally used to see if a circuit is performing as expected, but oscilloscopes are also useful for comparing different signals to each
other.
Oscilloscope Display: Example of an analog oscilloscope display. Shown is a Lissajous figure, showing a harmonic relationship of
one horizontal oscillation cycle to three vertical oscillation cycles.
Many oscilloscopes also use CRT displays, though LCD displays are becoming more common. In oscilloscope CRTs, electrostatic
deflection is used, rather than the magnetic deflection commonly used with television and other large CRTs. The beam is deflected
horizontally by applying an electric field between a pair of plates to its left and right, and vertically by applying an electric field to
plates above and below.
Oscilloscopes use electrostatic rather than magnetic deflection because the inductive reactance of the magnetic coils would limit
the frequency response of the instrument. The color of the oscilloscope phosphor is much less important than in the case of color
televisions or computer monitors since the primary purpose is to evaluate signal voltages rather than construct complex images;
however, the persistence of the phosphor may be more important. Phosphors are available with persistences ranging from less than
one microsecond to several seconds. For visual observation of brief transient events, a long persistence phosphor may be desirable.
For events which are fast and repetitive, or high frequency, a short-persistence phosphor is generally preferable.
Key Points
The primary components of a cathode ray tube (CRT) consist of a vacuum tube containing an electron gun and a screen lined
with phosphors. CRTs are used to produce images.
The phosphors in a CRT’s screen are the materials that directly produce the photons generated by the CRT. These phosphors are
struck by incoming electrons from the electron gun, absorb energy, and then re-emit some or all of that energy in the form of
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light.
CRT technology used to be common in televisions and computer monitors. Color CRTs contain three electron guns
corresponding to three types of phosphors, one for each primary color (red, blue, and green). Examples of monochromatic CRTs
include black and white TVs and old computer terminals.
Oscilloscopes, devices used to measure and display voltages, also use CRT displays. In this case, the persistence of the
phosphor is more important than the color.
The CRTs in televisions and computer monitors bend the electron beams with magnetic deflection, while oscilloscopes rely on
electrostatic deflection.
Key Terms
phosphor: A substance that exhibits the phenomenon of luminescence; often transition metal compounds or rare earth
compounds of various types. The most common uses of phosphors are in CRT displays and fluorescent lights.
raster: A scanning pattern of parallel lines that form the display of an image projected on a cathode-ray tube of a television set
or display screen.
electron gun: Any device that produces a stream of electrons, especially a narrow stream that is focused onto a phosphor
screen.
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CHAPTER OVERVIEW
This page titled 19: Electric Current and Resistance is shared under a not declared license and was authored, remixed, and/or curated by
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1
19.1: Overview
learning objectives
Explain difference between conductor and resistor
Electric Current
When we speak of electric current, often we are referring to a specific quantity – the rate at which charge flows. A large current,
such as that used to start a truck engine, moves a large amount of charge in a small time, whereas a small current, such as that used
to operate a hand-held calculator, moves a small amount of charge over a long period of time. In equation form, electric current I is
defined to be
ΔQ
I = (19.1.1)
ΔT
where Q is the amount of charge passing through a given area in time t. The SI unit for current is the ampere (A), named for the
French physicist André-Marie Ampère (1775–1836). Since I=ΔQ/Δt, we see that an ampere is one coulomb per second:
1A = 1C/s (19.1.2)
The flow of electricity requires a medium in which charge can flow. We call an object or medium that allows charge to flow a
conductor, while the empirical measure of a material’s ability to conduct charge is called the electrical conductance. The SI unit for
conductance is the siemens (S).
Electric Current: The rate of flow of charge is current. An ampere is the flow of one coulomb through an area in one second.
Resistance
The opposite of conductance is resistance – a quantity that describes how strongly a material opposes the flow of electric current.
An object or medium that has high electrical resistance is called a resistor. We will see that the resistance of an object depends on
its shape and the material of which it is composed. The SI unit for resistance is the ohm (symbol: ).
Electric Circuits
A useful and practical way to learn about electric current and resistance is to study circuits. The figure above shows a simple circuit
and the standard schematic representation of a battery, conducting path, and load (a resistor). Schematics are very useful in
visualizing the main features of a circuit. A single schematic can represent a wide variety of situations. The schematic in (b), for
example, can represent anything from a truck battery connected to a headlight lighting the street in front of the truck to a small
battery connected to a penlight lighting a keyhole in a door. Such schematics are useful because the analysis is the same for a wide
variety of situations. We need to understand a few schematics to apply the concepts and analysis to many more situations.
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Simple Electric Circuit: (a) A simple electric circuit. A closed path for current to flow through is supplied by conducting wires
connecting a load to the terminals of a battery. (b) In this schematic, the battery is represented by the two parallel red lines,
conducting wires are shown as straight lines, and the zigzag represents the load. The schematic represents a wide variety of similar
circuits.
Note that the direction of current flow in the figure is from positive to negative. The direction of conventional current is the
direction that positive charge would flow. Depending on the situation, positive charges, negative charges, or both may move. In
metal wires, for example, current is carried by electrons—that is, negative charges move. In ionic solutions, such as salt water, both
positive and negative charges move.
It is important to realize that there is an electric field in conductors responsible for producing the current. Unlike static electricity,
where a conductor in equilibrium cannot have an electric field in it, conductors carrying a current have an electric field and are not
in static equilibrium. An electric field is needed to supply energy to move the charges.
Armed with these basics, we’ll begin to tackle the harder details of this topic in the next section.
Key Points
Electric current is the movement of electric charge through a conductive medium.
We also use the term “current” as a quantity to describe the rate at which charge flows through a medium. The SI unit for
current is the ampere (A), which is equal to a coulomb per second (C/s).
Conductance is a quantity describing how easily charge can flow through a material, while resistance is the inverse, a measure
of how strongly a material opposes electric flow.
An object that allows charge to flow easily is called a conductor, while an object that resists the flow of charge is called a
resistor.
Key Terms
conductive medium: A material that can transmit electricity.
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electrical resistance: The opposition offered by an electrical conductor to the flow of a current through itself, resulting in a
conversion of electrical energy into heat and radiation. The SI derived unit of resistance is the ohm. Symbol: R.
electric charge: A quantum number that determines the electromagnetic interactions of some subatomic particles; by
convention, the electron has an electric charge of -1 and the proton +1, and quarks have fractional charge.
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19.2: Electric Current
learning objectives
Describe the functions and identify the major components of a battery
A battery is a device that converts chemical energy directly to electrical energy. It consists of a number of voltaic cells connected in
series by a conductive electrolyte containing anions and cations. One half-cell includes electrolyte and the anode, or negative
electrode; the other half-cell includes electrolyte and the cathode, or positive electrode. In the redox (reduction-oxidation) reaction
that powers the battery, cations are reduced (electrons are added) at the cathode, while anions are oxidized (electrons are removed)
at the anode. The electrodes do not touch each other but are electrically connected by the electrolyte. Some cells use two half-cells
with different electrolytes. A separator between half-cells allows ions to flow, but prevents mixing of the electrolytes.
Each half-cell has an electromotive force (or emf), determined by its ability to drive electric current from the interior to the exterior
of the cell. The net emf of the cell is the difference between the emfs of its half-cells, or the difference between the reduction
potentials of the half-reactions.
The electrical driving force across the terminals of a cell is known as the terminal voltage (difference) and is measured in volts.
When a battery is connected to a circuit, the electrons from the anode travel through the circuit toward the cathode in a direct
circuit. The voltage of a battery is synonymous with its electromotive force, or emf. This force is responsible for the flow of charge
through the circuit, known as the electric current.
A battery stores electrical potential from the chemical reaction. When it is connected to a circuit, that electric potential is converted
to kinetic energy as the electrons travel through the circuit. Electric potential is defined as the potential energy per unit charge (q).
The voltage, or potential difference, between two points is defined to be the change in potential energy of a charge q moved from
point 1 to point 2, divided by the charge. Rearranged, this mathematical relationship can be described as:
Voltage is not the same as energy. Voltage is the energy per unit charge. Thus a motorcycle battery and a car battery can both have
the same voltage (more precisely, the same potential difference between battery terminals), yet one stores much more energy than
the other. The car battery can move more charge than the motorcycle battery, although both are 12V batteries.
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Ideal and Real Batteries: A brief introduction to ideal and real batteries for students studying circuits.
Symbol of a Battery in a Circuit Diagram: This is the symbol for a battery in a circuit diagram. It originated as a schematic
drawing of the earliest type of battery, a voltaic pile. Notice the positive cathode and negative anode. This orientation is important
when drawing circuit diagrams to depict the correct flow of electrons.
learning objectives
Describe the relationship between the electrical current, voltage, and resistance in a circuit
To understand how to measure current and voltage in a circuit, you must also have a general understanding of how a circuit works
and how its electrical measurements are related.
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What is Voltage? Let's start here. Lesson 1 of 2
A Simple Circuit: A simple electric circuit made up of a voltage source and a resistor
According to Ohm’s law, The electrical current I, or movement of charge, that flows through most substances is directly
proportional to the voltage V applied to it. The electric property that impedes current (crudely similar to friction and air resistance)
is called resistance R. Collisions of moving charges with atoms and molecules in a substance transfer energy to the substance and
limit current. Resistance is inversely proportional to current. Ohm’s law can therefore be written as follows:
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V
I = (19.2.2)
R
where I is the current through the conductor in amperes, V is the potential difference measured across the conductor in volts, and R
is the resistance of the conductor in ohms (Ω). More specifically, Ohm’s law states that R in this relation is constant, independent of
the current. Using this equation, we can calculate the current, voltage, or resistance in a given circuit.
For example, if we had a 1.5V battery that was connected in a closed circuit to a lightbulb with a resistance of 5Ω, what is the
current flowing through the circuit? To solve this problem, we would just substitute the given values into Ohm’s law: I = 1.5V/5Ω;
I = 0.3 amperes. If we know the current and the resistance, we can rearrange the Ohm’s law equation and solve for voltage V:
V = IR (19.2.3)
learning objectives
Relate the drift velocity with the velocity of free charges in conductors
Drift Speed
Electrical signals are known to move very rapidly. Telephone conversations carried by currents in wires cover large distances
without noticeable delays. Lights come on as soon as a switch is flicked. Most electrical signals carried by currents travel at speeds
on the order of 108m/s, a significant fraction of the speed of light. Interestingly, the individual charges that make up the current
move much more slowly on average, typically drifting at speeds on the order of 10−4m/s.
The high speed of electrical signals results from the fact that the force between charges acts rapidly at a distance. Thus, when a free
charge is forced into a wire, the incoming charge pushes other charges ahead of it, which in turn push on charges farther down the
line. The resulting electrical shock wave moves through the system at nearly the speed of light. To be precise, this rapidly moving
signal or shock wave is a rapidly propagating change in the electric field.
Electrons Moving Through a Conductor: When charged particles are forced into this volume of a conductor, an equal number
are quickly forced to leave. The repulsion between like charges makes it difficult to increase the number of charges in a volume.
Thus, as one charge enters, another leaves almost immediately, carrying the signal rapidly forward.
Drift Velocity
Good conductors have large numbers of free charges in them. In metals, the free charges are free electrons. The distance that an
individual electron can move between collisions with atoms or other electrons is quite small. The electron paths thus appear nearly
random, like the motion of atoms in a gas. However, there is an electric field in the conductor that causes the electrons to drift in
the direction shown (opposite to the field, since they are negative). The drift velocity vdis the average velocity of the free charges
after applying the field. The drift velocity is quite small, since there are so many free charges. Given an estimate of the density of
free electrons in a conductor (the number of electrons per unit volume), it is possible to calculate the drift velocity for a given
current. The larger the density, the lower the velocity required for a given current.
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Drift Speed: Free electrons moving in a conductor make many collisions with other electrons and atoms. The path of one electron
is shown. The average velocity of the free charges is called the drift velocity and is in the direction opposite to the electric field for
electrons. The collisions normally transfer energy to the conductor, requiring a constant supply of energy to maintain a steady
current.
It is possible to obtain an expression for the relationship between the current and drift velocity by considering the number of free
charges in a segment of wire. The number of free charges per unit volume is given the symbol n and depends on the material. Ax is
the volume of a segment, so that the number of free charges in it is nAx. The charge ΔQ in this segment is thus qnAx, where q is the
amount of charge on each carrier. (Recall that for electrons, q is 1.60×10−19C. ) The current is the charge moved per unit time.
Thus, if all the original charges move out of this segment in time t, the current is:
ΔQ x
I = = qnA (19.2.4)
Δt Δt
Notably, x/Δt is the magnitude of the drift velocity vd, since the charges move an average distance x in a time t. Rearranging terms
gives: I = qnAvd, where I is the current through a wire of cross-sectional area A made of a material with a free charge density n. The
carriers of the current each have charges q and move with a drift velocity of magnitude vd.
Current density is the electric current per unit area of cross-section. It has units of Amperes per square meter.
Key Points
A battery stores electrical potential from the chemical reaction. When it is connected to a circuit, that electric potential is
converted to kinetic energy as the electrons travel through the circuit.
The voltage or potential difference between two points is defined to be the change in potential energy of a charge q moved from
point 1 to point 2, divided by the charge.
The voltage of a battery is synonymous with its electromotive force, or emf. This force is responsible for the flow of charge
through the circuit, known as the electric current.
A simple circuit consists of a voltage source and a resistor.
Ohm ‘s law gives the relationship between current I, voltage V, and resistance R in a simple circuit: I = V/R.
The SI unit for measuring the rate of flow of electric charge is the ampere, which is equal to a charge flowing through some
surface at the rate of one coulomb per second.
There is an electric field in conductors that causes electrons to drift in the direction opposite to the field. The drift velocity is the
average velocity of these free charges.
The expression for the relationship between the current and drift velocity can be obtained by considering the number of free
charges in a segment of wire.
I = qnAv relates the drift velocity to the current, where I is the current through a wire of cross-sectional area A made of a
material with a free charge density n. The carriers of the current each have a charge q and move with a drift velocity of
magnitude v.
Key Terms
battery: A device that produces electricity by a chemical reaction between two substances.
current: The time rate of flow of electric charge.
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voltage: The amount of electrostatic potential between two points in space.
electrical current: the movement of charge through a circuit
ohm: in the International System of Units, the derived unit of electrical resistance; the electrical resistance of a device across
which a potential difference of one volt causes a current of one ampere; symbol: Ω
ampere: A unit of electrical current; the standard base unit in the International System of Units. Abbreviation: amp. Symbol: A.
drift velocity: The average velocity of the free charges in a conductor.
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19.3: Resistance and Resistors
Contrast shape of current-voltage plots for ohmic and non-ohmic circuits
Ohm’s Law
What drives current? We can think of various devices—such as batteries, generators, wall outlets, and so on—which are necessary
to maintain a current. All such devices create a potential difference and are loosely referred to as voltage sources. When a voltage
source is connected to a conductor, it applies a potential difference V that creates an electric field. The electric field, in turn, exerts
force on charges, causing current. The current that flows through most substances is directly proportional to the voltage V applied
to it. The German physicist Georg Simon Ohm (1787-1854) was the first to experimentally demonstrate that the current in a metal
wire is directly proportional to the voltage applied:\ mathrmI ∝ V .
This important relationship is known as Ohm’s law. It can be viewed as a cause-and-effect relationship, with voltage the cause and
current the effect. This is an empirical law like that for friction—an experimentally observed phenomenon. Such a linear
relationship doesn’t always occur. Recall that while voltage drives current, resistance impedes it. Collisions of moving charges with
atoms and molecules in a substance transfer energy to the substance and limit current. The current is therefore inversely
proportional to the resistance: I ∝ .
R
1
Simple Circuit: A simple electric circuit in which a closed path for current to flow is supplied by conductors (usually metal wires)
connecting a load to the terminals of a battery, represented by the red parallel lines. The zigzag symbol represents the single resistor
and includes any resistance in the connections to the voltage source.
The unit for resistance is the ohm where 1Ω = 1 V/A. We can combine the two relations above to obtain I = V/R. This relationship
is also called Ohm’s law. In this form Ohm’s law really defines resistance for certain materials. Ohm’s law (like Hooke’s law) is not
universally valid. The many substances for which Ohm’s law holds are called ohmic. These include good conductors like copper
and aluminum, and some poor conductors under certain circumstances. Ohmic materials have a resistance R that is independent of
voltage V and current I. An object that has simple resistance is called a resistor, even if its resistance is small.
Voltage Drop: The voltage drop across a resistor in a simple circuit equals the voltage output of the battery.
Additional insight is gained by solving I=V/R for V, yielding V=IR. This expression for V can be interpreted as the voltage drop
across a resistor produced by the flow of current I. The phrase IR drop is often used for this voltage. If voltage is measured at
various points in a circuit, it will be seen to increase at the voltage source and decrease at the resistor. Voltage is similar to fluid
pressure. The voltage source is like a pump, creating a pressure difference, causing current—the flow of charge. The resistor is like
a pipe that reduces pressure and limits flow because of its resistance. Conservation of energy has important consequences here. The
voltage source supplies energy (causing an electric field and a current), and the resistor converts it to another form (such as thermal
energy). In a simple circuit (one with a single simple resistor), the voltage supplied by the source equals the voltage drop across the
19.3.1 https://phys.libretexts.org/@go/page/14559
resistor, since E=qΔV, and the same q flows through each. Thus, the energy supplied by the voltage source and the energy
converted by the resistor are equal.
In a true ohmic device, the same value of resistance will be calculated from R = V/I regardless of the value of the applied voltage
V. That is, the ratio of V/I is constant, and when current is plotted as a function of voltage the curve is linear (a straight line). If
voltage is forced to some value V, then that voltage V divided by measured current I will equal R. Or if the current is forced to
some value I, then the measured voltage V divided by that current I is also R. We visualize the plot of I versus V as a straight line.
There are, however, components of electrical circuits which do not obey Ohm’s law; that is, their relationship between current and
voltage (their I–V curve) is nonlinear (or non-ohmic). An example is the p-n junction diode.
Current-Voltage Curves: The I–V curves of four devices: two resistors, a diode, and a battery. The two resistors follow Ohm’s
law: The plot is a straight line through the origin. The other two devices do not follow Ohm’s law.
19.3.2 https://phys.libretexts.org/@go/page/14559
Ohm’s Law: A brief overview of Ohm’s Law.
learning objectives
Describe behaviors of a superconductor below a critical temperature and in a weak external magnetic field
Superconductivity is a phenomenon of exactly zero electrical resistance and expulsion of magnetic fields occurring in certain
materials when cooled below a characteristic critical temperature. It was discovered by Heike Kamerlingh Onnes (shown in ) on
April 8, 1911 in Leiden.
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Superconducting phase transition: Behavior of heat capacity (cv, blue) and resistivity (ρ, green) at the superconducting phase
transition.
Superconductors are also able to maintain a current with no applied voltage whatsoever—a property exploited in superconducting
electromagnets such as those found in MRI machines. Experiments have demonstrated that currents in superconducting coils can
persist for years without any measurable degradation. Experimental evidence points to a current lifetime of at least 100,000 years.
Theoretical estimates for the lifetime of a persistent current can exceed the estimated lifetime of the universe, depending on the
wire geometry and the temperature.
The value of this critical temperature varies from material to material. Usually, conventional superconductors have critical
temperatures ranging from around 20 K to less than 1 K. Solid mercury, for example, has a critical temperature of 4.2 K. As of
2009, the highest critical temperature found for a conventional superconductor is 39 K for magnesium diboride (MgB2), although
this material’s exotic properties cause some doubt about accurately classifying it as a “conventional” superconductor. High-
temperature superconductors can have much higher critical temperatures. For example, YBa2Cu3O7, one of the first cuprate
superconductors to be discovered, has a critical temperature of 92 K; mercury-based cuprates have been found with critical
temperatures in excess of 130 K. It is of note that the chemical composition and crystal structure of superconducting materials can
be quite complex, as seen in.
Unit Cell of YBaCuO superconductor: Unit Cell of YBaCuO superconductor. Atoms are indicated with different colors.
19.3.4 https://phys.libretexts.org/@go/page/14559
Resistance and Resistivity
Resistance and resistivity describe the extent to which an object or material impedes the flow of electric current.
learning objectives
Identify properties of the material that are described by the resistance and resistivity
Cylindrical Resistor: A uniform cylinder of length L and cross-sectional area A. Its resistance to the flow of current is similar to
the resistance posed by a pipe to fluid flow. The longer the cylinder, the greater its resistance. The larger its cross-sectional area A,
the smaller its resistance.
As mentioned, for a given shape, the resistance depends on the material of which the object is composed. Different materials offer
different resistance to the flow of charge. We define the resistivity ρ of a substance so that theresistanceR of an object is directly
proportional to ρ. Resistivity ρ is an intrinsic property of a material, independent of its shape or size. In contrast, the resistance R is
an extrinsic property that does depend on the size an shape of the resistor. (A similar intrinsic/extrinsic relation exists between heat
capacity C and the specific heat c). Recall that an object whose resistance is proportional to the voltage and current is known as a
resistor.
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Typical Resistor: A typical axial-lead resistor.
What determines resistivity? The resistivity of different materials varies by an enormous amount. For example, the conductivity of
teflon is about 1030 times lower than the conductivity of copper. Why is there such a difference? Loosely speaking, a metal has
large numbers of “delocalized” electrons that are not stuck in any one place, but free to move across large distances, whereas in an
insulator (like teflon), each electron is tightly bound to a single atom, and a great force is required to pull it away. Likewise,
resistors range over many orders of magnitude. Some ceramic insulators, such as those used to support power lines, have
resistances of 1012 Ω or more. A dry person may have a hand-to-foot resistance of 105 Ω, whereas the resistance of the human heart
is about 103 Ω. A meter-long piece of large-diameter copper wire may have a resistance of 10−5 Ω, and superconductors have no
resistance at all (they are non-ohmic). The potential difference (voltage) seen across the network is the sum of those voltages, thus
the total resistance (the series equivalent resistance) can be found as the sum of those resistances:
Req = R1 + R2 + ⋯ + RN (19.3.1)
As a special case, the resistance of N resistors connected in series, each of the same resistance R, is given by NR.Resistors in a
parallel configuration are each subject to the same potential difference (voltage), however the currents through them add. Thus the
equivalent resistance (Req) of the network can be computed:
1 1 1 1
= + +⋯ + (19.3.2)
Req R1 R2 RN
The parallel equivalent resistance can be represented in equations by two vertical lines “||” (as in geometry) as a simplified
notation. Occasionally two slashes “//” are used instead of “||”, in case the keyboard or font lacks the vertical line symbol. For the
case of two resistors in parallel, this can be calculated using:
R1 R2
Req = R1 ∥ R2 = (19.3.3)
R1 + R2
As a special case, the resistance of N resistors connected in parallel, each of the same resistance R, is given by R/N. A resistor
network that is a combination of parallel and series connections can be broken up into smaller parts that are either one or the other,
such as is shown in.
Resistor Network: In this combination circuit, the circuit can be broken up into a series component and a parallel component.
However, some complex networks of resistors cannot be resolved in this manner. These require a more sophisticated circuit
analysis. One practical application of these relationships is that a non-standard value of resistance can generally be synthesized by
connecting a number of standard values in series or parallel. This can also be used to obtain a resistance with a higher power rating
than that of the individual resistors used. In the special case of N identical resistors all connected in series or all connected in
parallel, the power rating of the individual resistors is thereby multiplied by N.
19.3.6 https://phys.libretexts.org/@go/page/14559
High School Physics - Resistance, Resistors, and Re…
Re…
Resistance, Resistors, and Resistivity: A brief overview of resistance, resistors, and resistivity.
learning objectives
Compare temperature dependence of resistivity and resistance for large and small temperature changes
The resistivity of all materials depends on temperature. Some materials can become superconductors (zero resistivity) at very low
temperatures (see ). Conversely, the resistivity of conductors increases with increasing temperature. Since the atoms vibrate more
rapidly and over larger distances at higher temperatures, the electrons moving through a metal, for example, create more collisions,
effectively making the resistivity higher. Over relatively small temperature changes (about 100ºC or less), resistivity ρ varies with
temperature change ΔT as expressed in the following equation:
19.3.7 https://phys.libretexts.org/@go/page/14559
Resistance of a sample of mercury: The resistance of a sample of mercury is zero at very low temperatures—it is a
superconductor up to about 4.2 K. Above that critical temperature, its resistance makes a sudden jump and then increases nearly
linearly with temperature.
p = p0 (1 + αΔT) (19.3.4)
where ρ0 is the original resistivity and α is the temperature coefficient of resistivity. For larger temperature changes, α may vary, or
a nonlinear equation may be needed to find ρ. For this reason it is usual to specify a suffix for the temperature at which the
substance was measured (such as α15) and the relationship only holds in a range of temperatures around the reference. Note that α is
positive for metals, meaning their resistivity increases with temperature. The temperature coefficient is typically +3×10−3 K−1 to
+6×10−3 K−1 for metals near room temperature. Some alloys have been developed specifically to have a small temperature
dependence. Manganin (made of copper, manganese and nickel), for example, has α close to zero, so its resistivity varies only
slightly with temperature. This is useful for making a temperature-independent resistance standard, for example.
Note also that α is negative for semiconductors, meaning that their resistivity decreases with increasing temperature. They become
better conductors at higher temperature because increased thermal agitation increases the number of free charges available to carry
current. This property of decreasing ρ with temperature is also related to the type and amount of impurities present in the
semiconductors.
The resistance of an object also depends on temperature, since R0 is directly proportional to ρ. For a cylinder we know R=ρL/A, so
if L and A do not change greatly with temperature, R will have the same temperature dependence as ρ. (Examination of the
coefficients of linear expansion shows them to be about two orders of magnitude less than typical temperature coefficients of
resistivity, and so the effect of temperature on L and A is about two orders of magnitude less than on ρ. ) Thus,
R = R0 (1 + αΔT) (19.3.5)
is the temperature dependence of the resistance of an object, where R0 is the original resistance and R is the resistance after a
temperature change T. Numerous thermometers are based on the effect of temperature on resistance (see ). One of the most
common is the thermistor, a semiconductor crystal with a strong temperature dependence, the resistance of which is measured to
obtain its temperature. The device is small so it quickly comes into thermal equilibrium with the part of a person it touches.
19.3.8 https://phys.libretexts.org/@go/page/14559
Thermometers: These familiar thermometers are based on the automated measurement of a thermistor’s temperature-dependent
resistance.
Key Points
Voltage drives current while resistance impedes it.
Ohm ‘s Law refers to the proportion relation between voltage and current. It also applies to the specific equation V=IR, which is
valid when considering circuits that contain simple resistors (whose resistance is independent of voltage and current).
Circuits or components that obey the relation V=IR are known as ohmic and have current-voltage plots that are linear and pass
through the origin.
There are components and circuits that are non-ohmic; their I-V plots are not linear and/or don’t pass through the origin.
Superconductivity is a Superconductivity is a thermodynamic phase and possesses certain distinguishing properties which are
largely independent of microscopic details.
In superconducting materials, the characteristics of superconductivity appear when the temperature is lowered below a critical
temperature. The onset of superconductivity is accompanied by abrupt changes in various physical properties.
When a superconductor is placed in a weak external magnetic field H, and cooled below its transition temperature, the magnetic
field is ejected.
Superconductors are able to maintain a current with no applied voltage.
The resistance of an object (i.e., a resistor ) depends on its shape and the material of which it is composed.
Resistivity ρ is an intrinsic property of a material and directly proportional to the total resistance R, an extrinsic quantity that
depends on the length and cross-sectional area of a resistor.
The resistivity of different materials varies by an enormous amount. Likewise, resistors range over many orders of magnitude.
Resistors are arranged in series or parallel configurations. The equivalent resistance of a network of resistors in series is the sum
of all the resistance. The inverse of the equivalent resistance of a network of resistors in parallel is the sum of the inverse of the
resistance of each resistor.
Over temperature changes of 100ºC or less, resistivity (ρ) varies with temperature change ΔT as: p=p0(1+αΔT)p=p0(1+αΔT)
where ρ0 is the original resistivity and α is the temperature coefficient of resistivity.
For large temperature changes, nonlinear variation of resistivity with temperature is observed.
The resistance of an object demonstrates similar temperature dependence as resistivity since resistance is directly proportional
to resistivity.
Key Terms
simple circuit: A circuit with a single voltage source and a single resistor.
ohmic: That which obeys Ohm’s law.
high-temperature superconductors: Materials that behave as superconductors at unusually high temperatures (above about 30
K).
critical temperature: In superconducting materials, the characteristics of superconductivity appears at (and continues below)
this temperature.
superconductivity: The property of a material whereby it has no resistance to the flow of an electric current.
series equivalent resistance: The resistance of a network of resistors arranged such that the voltage across the network is the
sum of the voltage across each resistor. In this case, the equivalent resistance is the sum of the resistance of all the resistors in
the network.
19.3.9 https://phys.libretexts.org/@go/page/14559
parallel equivalent resistance: the resistance of a network such that each resistor is subject to the same potential difference
(voltage), so that the currents through them add. In this case the inverse of the equivalent resistance is equal to the sum of the
inverse resistance of all the resistors in the network.
resistivity: In general, the resistance to electric current of a material; in particular, the degree to which a material resists the
flow of electricity.
temperature coefficient of resistivity: An empirical quantity, denoted by α, which describes the change in resistance or
resistivity of a material with temperature.
semiconductor: A substance with electrical properties intermediate between a good conductor and a good insulator.
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19.4: Electric Power and Energy
learning objectives
Formulate the relationship between the energy usage and the electric power
In many cases it is necessary to calculate the energy usage by an electric device or a collection of devices, such as in a home. For
example, we (or the electric power utility) may want to calculate the amount of money owed for electricity consumed. In another
case, we might need to determine the energy necessary to power a component or apparatus for a given period of time. The last
distinction is crucial – the energy used by a circuit or component is the time integral of the electric power.
Power
Recall that power is the rate at which work is done – or the rate at which energy is consumed or produced – and is measured in
watts (W). The electric power in watts produced by an electric current I consisting of a charge of Q coulombs every t seconds
QV
passing through an electric potential (voltage) difference of V is P = t
= IV , where Q is electric charge in coulombs, t is time
In resistive circuits where Ohm ‘s Law applies, the power can be expressed as P = I R =2 V
R
, where R is the electrical resistance.
Power isn’t necessarily constant; it may vary over time. The general expression for electric power is then
P(t) = I(t)V(t) (19.4.1)
Energy
In any given time interval, the energy consumed (or provided, depending on your perspective) is given by PE = qV , where E is
the electric energy, V is the voltage, and q is the amount of charge moved in the time interval under consideration. We can relate
the total energy consumed to the power by integrating over time:. Positive energy corresponds to consumed energy and negative
energy corresponds to energy production. Note that a circuit element having a power profile that is both positive and negative over
some time interval could consume or produce energy according to the sign of the integral of power. If the power is constant over
the time interval then the energy can be expressed simply as:
E = Pt (19.4.2)
19.4.1 https://phys.libretexts.org/@go/page/14560
Compact Fluorescent Light (CFL): CFLs are much more efficient than incandescent bulbs and so consume much less energy for
the intensity light produces.
Key Points
Recall that power is the rate work is done, or the rate at which energy is consumed or produced. In terms of current and voltage
it is P=IV.
The energy used is the amount of charge q moved through voltage V in a time interval t. It is equal to the integral of power over
time.
A common unit used to describe energy usage is the kilowatt-hour, the energy of 1000 W acting over one hour.
Key Terms
kilowatt-hour: a unit of electrical energy equal to that done by one kilowatt acting for one hour; equal to 3.6 mega-joules.
Symbol: kWh.
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19.5: Alternating Currents
learning objectives
Discuss applications of a phase vector
Phasors
Complex numbers play an important role in physics. Usually, complex numbers are written in terms of their real part plus the
imaginary part. For example, a + bi where a and b are real numbers, and ii signals the imaginary part. However, it is often practical
to write complex numbers in the form of an exponential called a phasor.
In physics, a phase vector, or phasor, is a representation of a sinusoidal function whose amplitude (A), frequency (ω), and phase (θ)
are time-invariant, as diagramed in. Phasors separate the dependencies on A , ω, and θ into three independent factors. This can be
particularly useful because the frequency factor (which includes the time-dependence of the sinusoid) is often common to all the
components of a linear combination of sinusoids. In those situations, phasors allow this common feature to be factored out, leaving
just the A and θ features. The result is that trigonometry reduces to algebra, and linear differential equations become algebraic ones.
The term phasor therefore often refers to just those two factors.
Phasor Diagram: An example of series RLC circuit and respective phasor diagram for a specific ω. Electrical engineers,
electronics engineers, electronic engineering technicians and aircraft engineers all use phasor diagrams to visualize complex
constants and variables (phasors). Like vectors, arrows drawn on graph paper or computer displays represent phasors.
Phasors are often used in electrical systems when considering voltages and currents that vary sinusoidally in time, such as in RLC
circuits.
Definition
Sinusoids can be represented mathematically as the sum of two complex-valued functions:
i(ωt+θ) −i(ωt+θ)
e +e
A ⋅ cos(ωt + θ) = A ⋅ (19.5.1)
2
As indicated above, phasor can refer to either Aeiθ⋅eiωtAeiθ⋅eiωt or just the complex constant, AeiθAeiθ. In the latter case, it is
understood to be a shorthand notation, encoding the amplitude and phase of an underlying sinusoid.
19.5.1 https://phys.libretexts.org/@go/page/14561
x(t) = Acos(ωt + θ) (19.5.3)
illustrated in the figure below is a cosinusoidal signal with amplitude A, frequency, and phase θ. The amplitude A characterizes the
peak-to-peak swing of 2A, the angular frequency ω characterizes the period T=2π/ω between negative- to-positive zero crossings
(or positive peaks or negative peaks), and the phase θ characterizes the time τ=−θ/ω when the signal reaches its first peak. With so
defined, the signal x(t) may also be written as
In this form the signal is easy to plot. Simply draw a cosinusoidal wave with amplitude A and period T; then strike the origin (t=0)
so that the signal reaches its peak at τ. In summary, the parameters that determine a cosinusoidal signal have the following units:
A, arbitrary (e.g., volts or meters/sec, depending upon the application )
ω, in radians /sec (rad/sec)
T, in seconds (sec)
θ, in radians (rad)
τ, in seconds (sec)
Sinusoidal Steady State and the Series RLC CircuitPhasors may be used to analyze the behavior of electrical and mechanical
systems that have reached a kind of equilibrium called sinusoidal steady state.
In the sinusoidal steady state, every voltage and current (or force and velocity) in a system is sinusoidal with angular frequency ω.
However, the amplitudes and phases of these sinusoidal voltages and currents are all different.
For example, the voltage across a resistor might lead the voltage across a capacitor by 90∘ and lag the voltage across an inductor by
90∘. In order to make our application of phasors to electrical systems concrete, we consider the series RLC circuit illustrated in.
The arrow labeled i(t) denotes a current that flows in response to the voltage applied.
19.5.2 https://phys.libretexts.org/@go/page/14561
Series RLC Circuit: Series RLC Circuit.
We will assume that the voltage source is an audio oscillator that produces the voltage:
V(t) = Acos(ωt + θ) (19.5.6)
We then describe the voltage source by the phasor V and remember that we can always compute the actual voltage by multiplying
by eiωtand taking the real part.
learning objectives
Relate the root mean square voltage and current in an alternating circut with the peak voltage and current and the average
power
Definition
The root mean square (abbreviated RMS or rms), also known as the quadratic mean, is a statistical measure of the magnitude of a
varying quantity. It is especially useful when the function alternates between positive and negative values, e.g., sinusoids.The RMS
value of a set of values (or a continuous-time function such as a sinusoid) is the square root of the arithmetic mean of the squares of
the original values (or the square of the function). In the case of a set of n values {x1,x2,….,xn}, the RMS value is given by this
formula:
−−−−−−−−−−−−−−−−− −
1
2 2 2
xrms =√ (x + x + ⋯ + xn ) (19.5.9)
1 2
n
The corresponding formula for a continuous function f(t)defined over the interval T1 ≤ t ≤ T2 is as follows:
−−−−−−−−−−−−−−−−−
T2
1
2
frms = √ ∫ [t(t)] dt (19.5.10)
T2 − T1 T1
19.5.3 https://phys.libretexts.org/@go/page/14561
The RMS for a function over all time is below.
−−−−−−−−−−−−
T
1 2
frms = lim √ ∫ [t(t)] dt (19.5.11)
T →∞ T 0
The RMS over all time of a periodic function is equal to the RMS of one period of the function. The RMS value of a continuous
function or signal can be approximated by taking the RMS of a series of equally spaced samples.
Sinusoidal Voltage and Current: (a) DC voltage and current are constant in time, once the current is established. (b) A graph of
voltage and current versus time for 60-Hz AC power. The voltage and current are sinusoidal and are in phase for a simple
resistance circuit. The frequencies and peak voltages of AC sources differ greatly.
V = V0 sin(2πf t) (19.5.12)
V is the voltage at time t, V0 is the peak voltage, and f is the frequency in hertz. For this simple resistance circuit, I=V/R, and so the
AC current is as follows:
I = I0 sin(2πf t) (19.5.13)
Here, I is the current at time t, and I0=V0/R is the peak current. Now using the definition above, let’s calculate the rms voltage and
rms current. First, we have
−−−−−−− T2
1 2
Vrms = √ ∫ [ V0 sin(ωt)] dt (19.5.14)
T2 − T1 T1
Here, we have replaced 2πf with ω. Since V0 is a constant, we can factor it out of the square root, and use a trig identity to replace
the squared sine function.
−−−−−−− T2
1 1 − cos(2ωt)
Vrms = V0 √ ∫ dt (19.5.15)
T2 − T1 T1
2
19.5.4 https://phys.libretexts.org/@go/page/14561
Since the interval is a whole number of complete cycles (per definition of RMS), the terms will cancel out, leaving:
−−−−−−−−−−−−− −−−−−−−−−−−−−−
1 t 1 T2 − T1
Vrm = V0 √ [ ] T1 = V0 √ (19.5.17)
T2 − T1 2 T2 − T1 2
V0
= – (19.5.18)
√2
Similarly, you can find that the RMS current can be expressed fairly simply:
–
Irms = I0 / √2 (19.5.19)
2
Vms
Pave = (19.5.22)
R
2
Pave = Irms R (19.5.23)
We can see from the above equations that we can express the average power as a function of the peak voltage and current (in the
case of sinusoidally varying current and voltage):
Average Power: AC power as a function of time. Since the voltage and current are in phase here, their product is non-negative and
fluctuates between zero and I0V0. Average power is (1/2)I0V0.
I0 V0 1
Pave = Irms Vrms = – – = V0 I0 (19.5.24)
√2 √2 2
The RMS values are also useful if the voltage varies by some waveform other than sinusoids, such as with a square, triangular or
sawtooth waves.
19.5.5 https://phys.libretexts.org/@go/page/14561
Waveforms: Sine, square, triangle, and sawtooth waveforms
learning objectives
Identify major risks associated with the electrical circuits and strategies to mitigate those risks
AC Circuit Lacking Safety Features: A schematic of a simple AC circuit with a voltage source and a single appliance represented
by the resistance R. It lacks safety features.
In practice, a simple AC circuit with no safety features is not how power is distributed. Modern household and industrial wiring
requires the three-wire system, which has several safety features. The first safety feature is the familiar circuit breaker (or fuse) that
prevents thermal overload. Secondly, there is a protective case around the appliance, as with a toaster or refrigerator. The case
prevents people from touching exposed wires and coming into electrical contact with the circuit, helping prevent shocks.
19.5.6 https://phys.libretexts.org/@go/page/14561
Three-Wire System: The three-wire system connects the neutral wire to the earth at the voltage source and the user location. It
exists at zero volts and supplies an alternative return path for the current through the earth. The case of the appliance is also
grounded to zero volts. A circuit breaker or fuse prevents thermal overload and exists in series on the active (live/hot) wire. Wire
insulation colors vary by region. It is essential to check locally to determine which color codes are in use, even if they were
followed in one particular installation.
There are three connections to the earth or ground (earth/ground, ). An earth/ground connection is a low-resistance path directly to
the earth. The two earth/ground connections on the neutral wire force it to exist at zero volts relative to the earth, giving the wire its
name. This wire is therefore safe to touch even if its insulation is missing. The neutral wire is the return path for the current to
follow in order to complete the circuit.
The two earth/ground connections supply an alternative path through the earth to complete the circuit, since the earth is a good
conductor. The earth/ground connection closest to the power source could be at the generating plant, while the other is situated at
the user’s location. The third earth/ground connection involves the case of the appliance, through the green earth/ground wire,
forcing the case to be at zero volts. The live or hot wire (live/hot) supplies the voltage and current to operate the appliance. The
three-wire system is connected to an appliance through a three-prong plug.
19.5.7 https://phys.libretexts.org/@go/page/14561
Three-Prong Plug: The standard three-prong plug can only be inserted one way to ensure the proper function of the three-wire
system.
Color-Coding
Insulating plastic is color-coded to identify live/hot, neutral, and ground wires, but these codes vary throughout the world. Live/hot
wires may be brown, red, black, blue, or grey. Neutral wires may be blue, black, or white. Since the same color may be used for
live/hot or neutral wires in different parts of the world, it is essential to confirm the color code for any given local region. The only
exception is the earth/ground wire, which is often green but may be yellow or ‘bare wire. ‘ Striped coatings are sometimes used for
the benefit of those who are colorblind.
Key Points
A phasor is a representation of a sinusoidal function whose amplitude (A), frequency (ω), and phase (θ) are time-invariant. If ω
is shared by all components of the system, it can be factored out, leaving just A and ω. The term phasor usually refers to the last
two factors.
Phasors greatly reduce the complexity of expressing sinusoidally varying signals.
Phasors may be used to analyze the behavior of electrical systems, such as RLC circuits, that have reached a kind of equilibrium
called sinusoidal steady state. In the sinusoidal steady state, every voltage and current in a system is sinusoidal with angular
frequency ω.
Phasors allow us to apply techniques used to solve DC circuits to solve RC circuits.
19.5.8 https://phys.libretexts.org/@go/page/14561
Recall that unlike DC current and voltage, which are constant, AC current and voltage vary over time. This is called alternating
current because the direction alternates.
The root mean square (abbreviated RMS or rms ) is a statistical measure of the magnitude of a varying quantity. We use the root
mean square to express the average current or voltage in an AC system.
The RMS current and voltage (for sinusoidal systems) are the peak current and voltage over the square root of two.
The average power in an AC circuit is the product of the RMS current and RMS voltage.
Electrical circuits carry the risks of overheating and potential electrical shocks.
Fuses and circuit breakers are used to stop currents that exceed a set safety limit, thus preventing overheating.
The three-wire system protects against thermal and shock hazards by using live, neutral, and ground wires, and grounding the
neutral wire and the conducting cases of appliances.
Before altering any circuitry, it is important to establish the correct color-coding scheme for your region (the color of live/hot,
neutral, and ground wires).
Alternating current has the potential to induce an EMF on the case of an appliance, which poses a shock hazard, so it is
important to ground the case.
Key Terms
sinusoidal steady state: Indicates every voltage and current in a system is sinusoidal with the same angular frequency ω.
complex numbers: Numbers that have an imaginary part. Usually represented as i.
phasor: A representation of a complex number in terms of a complex exponential.
root mean square: The square root of the arithmetic mean of the squares.
I0
rms current: the root mean square of the current, I rms= , where I0 is the peak current, in an AC system
√2
V0
rms voltage: the root mean square of the voltage, V rms =
√2
, where V0 is the peak voltage, in an AC system
thermal hazard: an electrical hazard caused by overheating (e.g., in a resistive element)
shock hazard: an electrical hazard that poses the risk of passing current through the body
three-wire system: a modern wiring system with safety precautions; contains live, neutral, and ground wires
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19.6: Electricity in the World
learning objectives
Identify factors that determine the lethality of an electric shock
There are two known categories of electrical hazards: thermal hazards and shock hazards. A thermal hazard is when excessive
electric power causes undesired thermal effects, such as starting a fire in the wall of a house. A shock hazard occurs when electric
current passes through a person. Shocks range in severity from painful but otherwise harmless to heart-stoppingly lethal.
Thermal Hazards
Electric power causes undesired heating effects whenever electric energy is converted to thermal energy at a rate faster than it can
be safely dissipated. A classic example of this is the short circuit, shown in. A short circuit is a low- resistance path between
terminals of a voltage source. Insulation on the appliance’s wires has worn through, allowing the two wires to come into contact.
Such an undesired contact with a high voltage is called a short. Since the resistance of the short, r, is very small, the power
dissipated in the short, P = V2/r, is very large. For example, if V is 120 V and r is 0.100 Ω, then the power is 144 kW, much greater
than that used by a typical household appliance. Thermal energy delivered at this rate will very quickly raise the temperature of
surrounding materials, melting or perhaps igniting them.
Short Circuit: A short circuit is an undesired low-resistance path across a voltage source. (a) Worn insulation on the wires of a
toaster allows them to come into contact with a low resistance r. Since P = V2/r, thermal power is created so rapidly that the cord
melts or burns. (b) A schematic of the short circuit.
A thermal hazard can be created even when a short circuit is not present if the wires in a circuit are overloaded with too much
current. The power dissipated in the supply wires is P = I2Rw, where Rw is the resistance of the wires and I the current flowing
through them. If either I or Rw is too large, the wires overheat. Thermal hazards can cause moderate to severe burns to those who
come in contact with the affected appliance or circuit.
Shock Hazards
Electric shock occurs upon contact of a body part with any source of electricity that causes a sufficient current through the skin,
muscles, or hair. Typically, the expression is used to describe an injurious exposure to electricity.The minimum current a human can
feel depends on the current type ( AC or DC ) and frequency. A person can feel at least 1 mA (rms) of AC current at 60 Hz and at
least 5 mA of DC current. The current may, if it is high enough, cause tissue damage or fibrillation, which leads to cardiac arrest.
60 mA of AC (rms, 60 Hz) or 300-500 mA of DC can cause fibrillation.The potential severity of the shock depends on paths
through the body that the currents take. If an electrical circuit is established by electrodes introduced in the body, bypassing the
skin, then the potential for lethality is much higher if a circuit through the heart is established. This is known as a microshock.
Currents of only 10 µA can be sufficient to cause fibrillation in this case.
A very dangerous possibility is the “can’t let go” effect illustrated in. The muscles that close the fingers are stronger than those that
open them, so the hand involuntarily closes around the wire shocking it. This can prolong the shock indefinitely. It can also be a
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danger to a person trying to rescue the victim, because the rescuer’s hand may close about the victim’s wrist. Usually the best way
to help the victim is to give the fist a hard blow with an insulator or to throw an insulator at the fist.
Electric Shock and Muscular Contractions: An electric current can cause muscular contractions with varying effects. (a) The
victim is “thrown” backward by involuntary muscle contractions that extend the legs and torso. (b) The victim can’t let go of the
wire that is stimulating all the muscles in the hand. Those that close the fingers are stronger than those that open them.
High Voltage Warning: International safety symbol “Caution, risk of electric shock” (ISO 3864), also known as the high-voltage
symbol
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learning objectives
Explain purpose of the electrocardiogram and identify functions performed by electric currents in the nerve system
Nerve Conduction
Electric currents in the complex system of nerves in our body allow us to sense the world, control parts of our body, and think.
There are three major functions of nerves. First, nerves carry messages from our sensory organs to the central nervous system,
consisting of the brain and spinal cord. Second, nerves carry messages from the central nervous system to muscles and other
organs. Third, nerves transmit and process signals within the central nervous system.
Nerve conduction is a general term for electrical signals carried by nerve cells. A voltage is created across the cell membrane of a
neuron in its resting state. This membrane separates electrically neutral fluids having differing concentrations of ions, the most
important varieties being Na+, K+, and Cl− (these are sodium, potassium, and chlorine ions). Free ions will diffuse from a region of
high concentration to one of low concentration. The cell membrane is semipermeable, meaning that some ions may cross it while
others cannot. In its resting state, the cell membrane is permeable to K+ and Cl−, and impermeable to Na+. Diffusion of K+ and Cl−
thus creates the layers of positive and negative charge on the outside and inside of the membrane, and the Coulomb force prevents
the ions from diffusing across in their entirety.
Creating a Voltage Across a Cell Membrane: The semipermeable membrane of a cell has different concentrations of ions inside
and out. Diffusion moves the K+ and Cl− ions in the direction shown, until the Coulomb force halts further transfer. This results in
a layer of positive charge on the outside, a layer of negative charge on the inside, and thus a voltage across the cell membrane. The
membrane is normally impermeable to Na+.
Once the charge layer has built up, the repulsion of like charges prevents more from moving across, and the attraction of unlike
charges prevents more from leaving either side. The result is two layers of charge right on the membrane, with diffusion being
balanced by the Coulomb force. A tiny fraction of the charges move across and the fluids remain neutral, while a separation of
charge and a voltage have been created across the membrane. Electric currents along the cell membrane are created by any stimulus
that changes the membrane’s permeability. The membrane thus temporarily becomes permeable to Na+, which then rushes in,
driven both by diffusion and the Coulomb force. This inrush of Na+ first neutralizes the inside membrane (called depolarization),
and then makes it slightly positive. The depolarization causes the membrane to again become impermeable to Na+, and the
movement of K+ quickly returns the cell to its resting potential, referred to as repolarization. This sequence of events results in a
voltage pulse, called the action potential and is shown in.
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Voltage channels are critical in the generation of an action potential: Top: view of an idealized action potential shows its
various phases as the action potential passes a point on a cell membrane. Bottom: Recordings of action potentials are often
distorted compared to the schematic view because of variations in electrophysiological techniques used to make the recording.
Only small fractions of the ions move, so that the cell can fire many hundreds of times without depleting the excess concentrations
of Na+ and K+. This is an example of active transport, wherein cell energy is used to move ions across membranes against
diffusion gradients and the Coulomb force. The action potential is a voltage pulse at one location on a cell membrane.
How does it get transmitted along the cell membrane as a nerve impulse ? The changing voltage and electric fields affect the
permeability of the adjacent cell membrane, so that the same process takes place there. The adjacent membrane depolarizes,
affecting the membrane farther down, and so on. Thus the action potential stimulated at one location triggers a nerve impulse that
moves slowly (about 1 m/s) along the cell membrane.
Electrocardiograms
Just as nerve impulses are transmitted by depolarization and repolarization of an adjacent membrane, the depolarization that causes
muscle contraction can also stimulate adjacent muscle cells to depolarize (fire) and contract. Thus, a depolarization wave can be
sent across the heart, coordinating its rhythmic contractions and enabling it to perform its vital function of propelling blood through
the circulatory system. An electrocardiogram (ECG) is a record of the voltages created by the wave of depolarization (and
subsequent repolarization) in the heart. Historically, ECGs were performed by placing electrodes on the left and right arms and the
left leg. The voltage between the right arm and the left leg is called the lead II potential and is an indicator of heart-muscle
function.
shows an ECG of the lead II potential and a graph of other major events during the cardiac cycle. The major features are labeled P,
Q, R, S, and T. The P wave is generated by the depolarization and contraction of the atria as they pump blood into the ventricles.
The QRS complex has a characteristic shape and time span, and is created by the depolarization of the ventricles as they pump
blood to the body. The lead II QRS signal also masks the repolarization of the atria. Finally, the T waveis generated by the
repolarization of the ventricles and is followed by the P wave in the next heartbeat. Arterial blood pressure varies with each part of
the heartbeat, with systolic (maximum) pressure occurring closely after the QRS complex, signaling the contraction of the
ventricles.
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ECG Curve: A normal ECG curve synchronized with other major events during the cardiac cycle.
learning objectives
Identify part of the heart that acts as a pacemaker
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Structure of the heart: Structure diagram of a coronal section of the human heart from an anterior view. The two larger chambers
are the ventricles.
The electric energy that stimulates the heart occurs in the sinoatrial node, which produces a definite potential and then discharges,
sending an impulse across the atria. In the atria the electrical signal moves from cell to cell (see section on nerve conduction and
the electrocardiogram ) while in the ventricles the signal is carried by specialized tissue called the Purkinje fibers which then
transmit the electric charge to the myocardium. shows the isolated heart conduction system.
Conduction System of the Heart: Isolated heart conduction system, showing the sinoatrial and Purkinje fibers.
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Sinoatrial Node Tissue: High magnification micrograph of sinoatrial node tissue and an adjacent nerve fiber.
Cells in the SA node, located in the upper right corner of the heart, will typically discharge (create action potentials) at about 60-
100 beats/minute. Because the sinoatrial node is responsible for the rest of the heart’s electrical activity, it is sometimes called the
primary pacemaker. If the SA node does not function, or the impulse generated in the SA node is blocked before it travels down the
electrical conduction system, a group of cells further down the heart will become the heart’s pacemaker. These cells form the
atrioventricular node (AV node), which is an area between the atria and ventricles, within the atrial septum. If the AV node also
fails, Purkinje fibers are capable of acting as the pacemaker. The reason Purkinje cells do not normally control the heart rate is that
they generate action potentials at a lower frequency than the AV or SA nodes.
Purkinje Fibers
The Purkinje fibers are located in the inner ventricular walls of the heart. These fibers consist of specialized cardiomyocytes that
are able to conduct cardiac action potentials more quickly and efficiently than any other cells in the heart. Purkinje fibers allow the
heart’s conduction system to create synchronized contractions of its ventricles, and are therefore essential for maintaining a
consistent heart rhythm.
During the ventricular contraction portion of the cardiac cycle, the Purkinje fibers carry the contraction impulse from both the left
and right bundle branch to the myocardium of the ventricles. This causes the muscle tissue of the ventricles to contract, thus
enabling a force to eject blood out of the heart. Atrial and ventricular discharge through the Purkinje trees is assigned on a standard
Electrocardiogram as the P Wave and QRS complex, respectively.
Purkinje fibers also have the ability of automatically firing at a rate of 15-40 beats per minute if left to their own devices. In
contrast, the SA node outside of parasympathetic control can fire a rate of almost 100 beats per minute. In short, they generate
action potentials, but at a slower rate than sinoatrial node and other atrial ectopic pacemakers. Thus they serve as the last resort
when other pacemakers fail.
Key Points
Electric power causes undesired heating effects whenever electric energy is converted to thermal energy at a rate faster than it
can be safely dissipated, such as in the case of a short circuit. This is referred to as a thermal hazard.
A shock hazard occurs when electric current passes through a person. If the current is high enough it can cause tissue damage or
fibrillation, which can lead to cardiac arrest.
The lethality of an electric shock is dependent on many variables, including current, duration, pathway, and the presence or
absence of very high voltage.
Nerves carry messages from our sensory organs to the central nervous system, from the central nervous system to muscles, and
within the central nervous system itself.
The effects of diffusion and the Coulomb force act together to allow ions like Na+, K+, and Cl- to create a voltage across cell
membranes.
An action potential is an event which alters the permeability of a cell membrane due to electric currents.
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The depolarization that causes muscle contraction can also stimulate contraction in adjacent muscle cells. A depolarization
wave across the heart is responsible for rhythmic contractions. An electrocardiogram (ECG) is a record of the voltages in the
heart.
The human heart provides continuous blood circulation through the cardiac cycle and is unsurprisingly one of the most vital
organs in the human body.
The heart is divided into four main chambers: the two upper chambers are called the left and right atria (singular atrium) and
two lower chambers are called the right and left ventricles. See for an illustration. There is a conduction system that transports
impulses through the heart.
The electric energy that stimulates the heart occurs in the sinoatrial node, which produces a definite potential and then
discharges, sending an impulse across the atria. It is thus the generator of normal sinus rhythm and functions as the heart’s
pacemaker.
In the atria the electrical signal moves from cell to cell while in the ventricles the signal is carried by specialized tissue called
the Purkinje fibers.
Purkinje fibers allow the heart’s conduction system to create synchronized contractions of its ventricles, and are therefore
essential for maintaining a consistent heart rhythm.
Key Terms
shock hazard: an electrical hazard that poses the risk of passing current through the body
thermal hazard: an electrical hazard caused by overheating (e.g., in a resistive element)
fibrillation: the rapid, irregular, and unsynchronized contraction of the muscle fibers of the heart
diffusion: the intermingling of the molecules of a fluid due to random thermal agitation
electrocardiogram: The visual output that an electrocardiograph produces
action potential: A short term change in the electrical potential that travels along a cell such as a nerve or muscle fiber.
sinoatrial node: a group of specialized cardiac muscle cells (tissue) located in the right atrium of the heart, which generates the
impulses that establish the normal sinus rhythm.
Purkinje fibers: specialized cardiac muscle cells that are able to conduct cardiac muscle potentials quickly and efficiently;
essential for maintaining consistent heart rhythm.
myocardium: The middle of the three layers forming the wall of the heart.
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CHAPTER OVERVIEW
2: Kinematics
Topic hierarchy
2.1: Basics of Kinematics
2.2: Speed and Velocity
2.3: Acceleration
2.4: Problem-Solving for Basic Kinematics
2.5: Free-Falling Objects
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1
2.1: Basics of Kinematics
Defining Kinematics
Kinematics is the study of the motion of points, objects, and groups of objects without considering the causes of its motion.
learning objectives
Define kinematics
Kinematics is the branch of classical mechanics that describes the motion of points, objects and systems of groups of objects,
without reference to the causes of motion (i.e., forces ). The study of kinematics is often referred to as the “geometry of motion.”
Objects are in motion all around us. Everything from a tennis match to a space-probe flyby of the planet Neptune involves motion.
When you are resting, your heart moves blood through your veins. Even in inanimate objects there is continuous motion in the
vibrations of atoms and molecules. Interesting questions about motion can arise: how long will it take for a space probe to travel to
Mars? Where will a football land if thrown at a certain angle? An understanding of motion, however, is also key to understanding
other concepts in physics. An understanding of acceleration, for example, is crucial to the study of force.
To describe motion, kinematics studies the trajectories of points, lines and other geometric objects, as well as their differential
properties (such as velocity and acceleration). Kinematics is used in astrophysics to describe the motion of celestial bodies and
systems; and in mechanical engineering, robotics and biomechanics to describe the motion of systems composed of joined parts
(such as an engine, a robotic arm, or the skeleton of the human body).
A formal study of physics begins with kinematics. The word “kinematics” comes from a Greek word “kinesis” meaning motion,
and is related to other English words such as “cinema” (movies) and “kinesiology” (the study of human motion). Kinematic
analysis is the process of measuring the kinematic quantities used to describe motion. The study of kinematics can be abstracted
into purely mathematical expressions, which can be used to calculate various aspects of motion such as velocity, acceleration,
displacement, time, and trajectory.
Kinematics of a particle trajectory: Kinematic equations can be used to calculate the trajectory of particles or objects. The
physical quantities relevant to the motion of a particle include: mass m, position r, velocity v, acceleration a.
learning objectives
Evaluate displacement within a frame of reference.
In order to describe the motion of an object, you must first describe its position — where it is at any particular time. More
precisely, you need to specify its position relative to a convenient reference frame. Earth is often used as a reference frame, and we
often describe the position of objects related to its position to or from Earth. Mathematically, the position of an object is generally
represented by the variable x.
Frames of Reference
There are two choices you have to make in order to define a position variable x. You have to decide where to put x = 0 and which
direction will be positive. This is referred to as choosing a coordinate system, or choosing a frame of reference. As long as you are
consistent, any frame is equally valid. But you don’t want to change coordinate systems in the middle of a calculation. Imagine
sitting in a train in a station when suddenly you notice that the station is moving backward. Most people would say that they just
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failed to notice that the train was moving — it only seemed like the station was moving. But this shows that there is a third
arbitrary choice that goes into choosing a coordinate system: valid frames of reference can differ from each other by moving
relative to one another. It might seem strange to use a coordinate system moving relative to the earth — but, for instance, the frame
of reference moving along with a train might be far more convenient for describing things happening inside the train. Frames of
reference are particularly important when describing an object’s displacement.
FRAMES OF REFERENCE by Professor Hume and Professor Donald Ivey of the University of Toronto
In this classic film, Professors Hume and Ivey cleverly illustrate reference frames and distinguish between fixed and moving
frames of reference.
Frames of Reference (1960) Educational Film: Frames of Reference is a 1960 educational film by Physical Sciences Study
Committee. The film was made to be shown in high school physics courses. In the film University of Toronto physics professors
Patterson Hume and Donald Ivey explain the distinction between inertial and nonintertial frames of reference, while demonstrating
these concepts through humorous camera tricks. For example, the film opens with Dr. Hume, who appears to be upside down,
accusing Dr. Ivey of being upside down. Only when the pair flip a coin does it become obvious that Dr. Ivey — and the camera —
are indeed inverted. The film’s humor serves both to hold students’ interest and to demonstrate the concepts being discussed. This
PSSC film utilizes a fascinating set consisting of a rotating table and furniture occupying surprisingly unpredictable spots within
the viewing area. The fine cinematography by Abraham Morochnik, and funny narration by University of Toronto professors
Donald Ivey and Patterson Hume is a wonderful example of the fun a creative team of filmmakers can have with a subject that
other, less imaginative types might find pedestrian. Producer: Richard Leacock Production Company: Educational Development
Corp. Sponsor: Eric Prestamon
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Displacement
Displacement is the change in position of an object relative to its reference frame. For example, if a car moves from a house to a
grocery store, its displacement is the relative distance of the grocery store to the reference frame, or the house. The word
“displacement” implies that an object has moved or has been displaced. Displacement is the change in position of an object and can
be represented mathematically as follows:
Δx = xf − x0 (2.1.1)
shows the importance of using a frame of reference when describing the displacement of a passenger on an airplane.
Displacement in Terms of Frame of Reference: A passenger moves from his seat to the back of the plane. His location relative to
the airplane is given by x. The -4.0m displacement of the passenger relative to the plane is represented by an arrow toward the rear
of the plane. Notice that the arrow representing his displacement is twice as long as the arrow representing the displacement of the
professor (he moves twice as far).
learning objectives
Distinguish the difference between scalars and vectors
What is the difference between distance and displacement? Whereas displacement is defined by both direction and magnitude,
distance is defined by magnitude alone. Displacement is an example of a vector quantity. Distance is an example of a scalar
quantity. A vector is any quantity with both magnitude and direction. Other examples of vectors include a velocity of 90 km/h east
and a force of 500 newtons straight down.
2.1.3 https://phys.libretexts.org/@go/page/14438
Scalars and Vectors
Scalars and Vectors: Mr. Andersen explains the differences between scalar and vectors quantities. He also uses a demonstration to
show the importance of vectors and vector addition.
In mathematics, physics, and engineering, a vector is a geometric object that has a magnitude (or length) and direction and can be
added to other vectors according to vector algebra. The direction of a vector in one-dimensional motion is given simply by a plus
(+) or minus (−) sign. A vector is frequently represented by a line segment with a definite direction, or graphically as an arrow,
connecting an initial point A with a terminal point B, as shown in.
Vector representation: A vector is frequently represented by a line segment with a definite direction, or graphically as an arrow,
connecting an initial point A with a terminal point B.
Some physical quantities, like distance, either have no direction or no specified direction. In physics, a scalar is a simple physical
quantity that is not changed by coordinate system rotations or translations. It is any quantity that can be expressed by a single
number and has a magnitude, but no direction. For example, a 20ºC temperature, the 250 kilocalories (250 Calories) of energy in a
candy bar, a 90 km/h speed limit, a person’s 1.8 m height, and a distance of 2.0 m are all scalars, or quantities with no specified
direction. Note, however, that a scalar can be negative, such as a −20ºC temperature. In this case, the minus sign indicates a point
on a scale rather than a direction. Scalars are never represented by arrows. (A comparison of scalars vs. vectors is shown in. )
2.1.4 https://phys.libretexts.org/@go/page/14438
Scalars vs. Vectors: A brief list of quantities that are either scalars or vectors.
Key Points
To describe motion, kinematics studies the trajectories of points, lines and other geometric objects.
The study of kinematics can be abstracted into purely mathematical expressions.
Kinematic equations can be used to calculate various aspects of motion such as velocity, acceleration, displacement, and time.
Choosing a frame of reference requires deciding where the object’s initial position is and which direction will be considered
positive.
Valid frames of reference can differ from each other by moving relative to one another.
Frames of reference are particularly important when describing an object’s displacement.
Displacement is the change in position of an object relative to its reference frame.
A vector is any quantity that has magnitude and direction.
A scalar is any quantity that has magnitude but no direction.
Displacement and velocity are vectors, whereas distance and speed are scalars.
Key Terms
kinematics: The branch of mechanics concerned with objects in motion, but not with the forces involved.
displacement: A vector quantity that denotes distance with a directional component.
frame of reference: A coordinate system or set of axes within which to measure the position, orientation, and other properties
of objects in it.
scalar: A quantity that has magnitude but not direction; compare vector.
vector: A directed quantity, one with both magnitude and direction; the between two points.
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2.2: Speed and Velocity
Average Velocity: A Graphical Interpretation
Average velocity is defined as the change in position (or displacement) over the time of travel.
learning objectives
Contrast speed and velocity in physics
In everyday usage, the terms “speed” and “velocity” are used interchangeably. In physics, however, they are distinct quantities.
Speed is a scalar quantity and has only magnitude. Velocity, on the other hand, is a vector quantity and so has both magnitude and
direction. This distinction becomes more apparent when we calculate average speed and velocity.
Average speed is calculated as the distance traveled over the total time of travel. In contrast, average velocity is defined as the
change in position (or displacement) over the total time of travel.
Average Velocity: The kinematic formula for calculating average velocity is the change in position over the time of travel.
The SI unit for velocity is meters per second, or m/s, but many other units (such as km/h, mph, and cm/s) are commonly used.
Suppose, for example, an airplane passenger took five seconds to move -4 m (the negative sign indicates that displacement is
toward the back of the plane ). His average velocity would be:
Δx −4 m m
v = = = −0.8 (2.2.1)
t 5 s s
The minus sign indicates that the average velocity is also toward the rear of the plane.
The average velocity of an object does not tell us anything about what happens to it between the starting point and ending point,
however. For example, we cannot tell from average velocity whether the airplane passenger stops momentarily or backs up before
he gets to the back of the plane. To get more details, we must consider smaller segments of the trip over smaller time intervals.
To illustrate the difference between average speed and average velocity, consider the following additional example. Imagine you
are walking in a small rectangle. You walk three meters north, four meters east, three meters south, and another four meters west.
The entire walk takes you 30 seconds. If you are calculating average speed, you would calculate the entire distance (3 + 4 + 3 + 4 =
14 meters) over the total time, 30 seconds. From this, you would get an average speed of 14/30 = 0.47 m/s. When calculating
average velocity, however, you are looking at the displacement over time. Because you walked in a full rectangle and ended up
exactly where you started, your displacement is 0 meters. Therefore, your average velocity, or displacement over time, would be 0
m/s.
Average Speed vs. Average Velocity: If you started walking from one corner and went all the way around the rectangle in 30
seconds, your average speed would be 0.47 m/s, but your average velocity would be 0 m/s.
2.2.1 https://phys.libretexts.org/@go/page/14439
Instananeous Velocity: A Graphical Interpretation
Instantaneous velocity is the velocity of an object at a single point in time and space as calculated by the slope of the tangent line.
learning objectives
Differentiate instantaneous velocity from other ways of determining velocity
Typically, motion is not with constant velocity nor speed. While driving in a car, for example, we continuously speed up and slow
down. A graphical representation of our motion in terms of distance vs. time, therefore, would be more variable or “curvy” rather
than a straight line, indicating motion with a constant velocity as shown below. (We limit our discussion to one dimensional
motion. It should be straightforward to generalize to three dimensional cases.)
Motion with Changing Velocity: Motion is often observed with changing velocity. This would result in a curvy line when graphed
with distance over time.
To calculate the speed of an object from a graph representing constant velocity, all that is needed is to find the slope of the line; this
would indicate the change in distance over the change in time. However, changing velocity it is not as straightforward.
Since our velocity is constantly changing, we can estimate velocity in different ways. One way is to look at our instantaneous
velocity, represented by one point on our curvy line of motion graphed with distance vs. time. In order to determine our velocity at
any given moment, we must determine the slope at that point. To do this, we find a line that represents our velocity in that moment,
shown graphically in. That line would be the line tangent to the curve at that point. If we extend this line, we can easily calculate
the displacement of distance over time and determine our velocity at that given point. The velocity of an object at any given
moment is the slope of the tangent line through the relevant point on its x vs. t graph.
Determining instantaneous velocity: The velocity at any given moment is dened as the slope of the tangent line through the
relevant point on the graph
2.2.2 https://phys.libretexts.org/@go/page/14439
Instantaneous Velocity, Acceleration, Jerk, Slopes, …
Instantaneous Velocity, Acceleration, Jerk, Slopes, Graphs vs. Time: This is how kinematics begins.
In calculus, finding the slope of curve f (x) at x = x x=x0 is equivalent to finding the first derivative:
0
df (x)
|x=x . (2.2.2)
0
dx
One interpretation of this definition is that the velocity shows how many meters the object would travel in one second if it
continues moving at the same speed for at least one second.
Key Points
Average velocity can be calculated by determining the total displacement divided by the total time of travel.
The average velocity of an object does not tell us anything about what happens to it between the starting point and ending point.
Average velocity is different from average speed in that it considers the direction of travel and the overall change in position.
When velocity is constantly changing, we can estimate our velocity by looking at instantaneous velocity.
Instantaneous velocity is calculated by determining the slope of the line tangent to the curve at the point of interest.
Instantaneous velocity is similar to determining how many meters the object would travel in one second at a specific moment.
Key Terms
velocity: A vector quantity that denotes the rate of change of position with respect to time, or a speed with a directional
component.
instantaneous: (As in velocity)—occurring, arising, or functioning without any delay; happening within an imperceptibly brief
period of time.
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2.3: Acceleration
Graphical Interpretation
The graphical representation of acceleration over time can be derived through the graph of an object’s position over time.
learning objectives
Distinguish the difference between how to plot a velocity graph and how to plot an acceleration graph
In physics, acceleration is the rate at which the velocity of a body changes with time. It is a vector quantity with both magnitude
and direction. Acceleration is accompanied by a force, as described by Newton’s Second Law; the force, as a vector, is the product
of the mass of the object being accelerated and the acceleration (vector), or F = ma . The SI unit of acceleration is the meter per
second squared: m
s
2
Acceleration is a vector that points in the same direction as the change in velocity, though it may not always be in the direction of
motion. For example, when an object slows down, or decelerating, its acceleration is in the opposite direction of its motion.
The motion of an object can be depicted graphically by plotting the position of an object over time. This distance-time graph can be
used to create another graph that shows changes in velocity over time. Because acceleration is velocity in m
s
divided by time in s,
we can further derive a graph of acceleration from a graph of an object’s speed or position.
is a graph of an object’s position over time. This graph is similar to the motion of a car. In the beginning, the object’s position
changes slowly as it gains speed. In the middle, the speed is constant and the position changes at a constant rate. As it slows down
toward the end, the position changes more slowly. From this graph, we can derive a velocity vs time graph.
Position vs Time Graph: Notice that the object’s position changes slowly at the beginning of the journey, then more and more
quickly as it picks up speed. Its position then changes more slowly as it slows down at the end of the journey. In the middle of the
journey, while the velocity remains constant, the position changes at a constant rate.
This shows the velocity of the object over time. The object’s velocity increases in the beginning as it accelerates at the beginning,
then remains constant in the middle before it slows down toward the end. Notice that this graph is a representation of the slope of
the previous position vs time graph. From this graph, we can further derive an acceleration vs time graph.
2.3.1 https://phys.libretexts.org/@go/page/14440
Velocity vs Time: The object’s velocity increases as it accelerates at the beginning of the journey. It remains the same in the middle
of the journey (where there is no acceleration). It decreases as the object decelerates at the end of the journey.
To do this, we would also plot the slope of the velocity vs time graph. In this graph, the acceleration is constant in the three
different stages of motion. As we noted earlier, the object is increasing speed and changing positions slowly in the beginning. The
acceleration graph shows that the object was increasing at a positive constant acceleration during this time. In the middle, when the
object was changing position at a constant velocity, the acceleration was 0. This is because the object is no longer changing its
velocity and is moving at a constant rate. Towards the end of the motion, the object slows down. This is depicted as a negative
value on the acceleration graph. Note that in this example, the motion of the object is still forward (positive), but since it is
decelerating, the acceleration is negative.
Acceleration vs Time Graph: The object has positive acceleration as it speeds up at the beginning of the journey. It has no
acceleration as it travels at constant velocity in the middle of the journey. Its acceleration is negative as it slows down at the end of
the journey.
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High School Physics: Graphing Motion
learning objectives
Describe how constant acceleration affects the motion of an object
2.3.3 https://phys.libretexts.org/@go/page/14440
One-Dimensional Motion: When you drop an object, it falls vertically toward the center of the earth due to the constant
acceleration of gravity.
An object experiencing constant acceleration has a velocity that increases or decreases by an equal amount for any constant period
of time. Acceleration can be derived easily from basic kinematic principles. It is defined as the first time derivative of velocity (so
the second derivative of position with respect to time):
2
∂v ∂ x
a = = (2.3.1)
2
∂t ∂t
Assuming acceleration to be constant does not seriously limit the situations we can study and does not degrade the accuracy of our
treatment, because in a great number of situations, acceleration is constant. When it is not, we can either consider it in separate
parts of constant acceleration or use an average acceleration over a period of time.
The motion of falling objects is a simple, one-dimensional type of projectile motion in which there is no horizontal movement. For
example, if you held a rock out and dropped it, the rock would fall only vertically downward toward the earth. If you were to throw
the rock instead of just dropping it, it would follow a more projectile-like pattern, similar to the one a kicked ball follows.
Projectile motion is the motion of an object thrown or projected into the air and is subject only to the acceleration of gravity. The
object thrown is called a projectile, and the object’s path is called its trajectory. In two-dimensional projectile motion, there is both
a vertical and a horizontal component.
Due to the algebraic properties of constant acceleration, there are kinematic equations that relate displacement, initial velocity, final
velocity, acceleration, and time. A summary of these equations is given below.
x = x0 + v̄t (2.3.2)
v0 + v
v̄ = (2.3.3)
2
v = v0 + at (2.3.4)
1
2
x = x0 + v0 t + at (2.3.5)
2
2 2
v =v + 2a(x − x0 ) (2.3.6)
0
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Constant Acceleration Explained with Vectors and A…
A…
Constant Acceleration Explained with Vectors and Algebra: This video answers the question “what is acceleration? ".
Key Points
Acceleration is the rate at which the velocity of a body changes with time.
Acceleration is a vector that points in the same direction as the change in velocity, though it may not always be in the direction
of motion.
Because acceleration is velocity in m/s divided by time in s, we can derive a graph of acceleration from a graph of an object’s
speed or position.
Assuming acceleration to be constant does not seriously limit the situations we can study and does not degrade the accuracy of
our treatment.
Due to the algebraic properties of constant acceleration, there are kinematic equations that can be used to calculate
displacement, velocity, acceleration, and time.
Calculations with constant acceleration can be done in relation to one-dimensional motion as well as two-dimensional motion.
2.3.5 https://phys.libretexts.org/@go/page/14440
Key Terms
acceleration: The amount by which a speed or velocity increases (and so a scalar quantity or a vector quantity).
velocity: A vector quantity that denotes the rate of change of position with respect to time, or a speed with a directional
component.
position: A place or location.
kinematic: of or relating to motion or kinematics
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2.4: Problem-Solving for Basic Kinematics
Applications
There are four kinematic equations that describe the motion of objects without consideration of its causes.
learning objectives
Choose which kinematics equation to use in problems in which the initial starting position is equal to zero
Kinematics is the branch of classical mechanics that describes the motion of points, bodies (objects), and systems of bodies (groups
of objects) without consideration of the causes of motion. There are four kinematic equations when the initial starting position is
the origin, and the acceleration is constant:
1. v = v + at
0
2. d = (v + v)t or alternatively v
1
2
0 average =
d
t
2
3. d = v t + ( )
0
at
4. v = v + 2ad
2 2
0
Notice that the four kinematic equations involve five kinematic variables: d, v, v , a and t . Each of these equations contains only
0
four of the five variables and has a different one missing. This tells us that we need the values of three variables to obtain the value
of the fourth and we need to choose the equation that contains the three known variables and one unknown variable for each
specific situation.
Here the basic problem solving steps to use these equations:
Step one – Identify exactly what needs to be determined in the problem (identify the unknowns).
Step two – Find an equation or set of equations that can help you solve the problem.
Step three – Substitute the knowns along with their units into the appropriate equation, and obtain numerical solutions complete
with units.
Step four – Check the answer to see if it is reasonable: Does it make sense?
Problem-solving skills are obviously essential to success in a quantitative course in physics. More importantly, the ability to apply
broad physical principles, usually represented by equations, to specific situations is a very powerful form of knowledge. It is much
more powerful than memorizing a list of facts. Analytical skills and problem-solving abilities can be applied to new situations,
whereas a list of facts cannot be made long enough to contain every possible circumstance. Such analytical skills are useful both for
solving problems in a physics class and for applying physics in everyday and professional life.
Motion Diagrams
A motion diagram is a pictorial description of an object’s motion and represents the position of an object at equally spaced time
intervals.
learning objectives
Construct a motion diagram
A motion diagram is a pictorial description of the motion of an object. It displays the object’s location at various equally spaced
times on the same diagram; shows an object’s initial position and velocity; and presents several spots in the center of the diagram.
These spots reveal whether or not the object has accelerated or decelerated. For simplicity, the object is represented by a simple
shape, such as a filled circle, which contains information about an object’s position at particular time instances. For this reason, a
motion diagram is more information than a path diagram. It may also display the forces acting on the object at each time instance.
is a motion diagram of a simple trajectory. Imagine the object as a hockey puck sliding on ice. Notice that the puck covers the same
distance per unit interval along the trajectory. We can conclude that the puck is moving at a constant velocity and, therefore, there is
no acceleration or deceleration during the motion.
2.4.1 https://phys.libretexts.org/@go/page/14441
Puck Sliding on Ice: Motion diagram of a puck sliding on ice. The puck is moving at a constant velocity.
One major use of motion diagrams is the presentation of film through a series of frames taken by a camera; this is sometimes called
stroboscopic technique (as seen in ). Viewing an object on a motion diagram allows one to determine whether an object is speeding
up or slowing down, or if it is at constant rest. As the frames are taken, we can assume that an object is at a constant rest if it
occupies the same position over time. We can assume that an object is speeding up if there is a visible increase in the space between
objects as time passes, and that it is slowing down if there is a visible decrease in the space between objects as time passes. The
objects on the frame come very close together.
Bouncing Ball: A bouncing ball captured with a stroboscopic flash at 25 images per second.
Key Points
The four kinematic equations involve five kinematic variables: d, v, v , a and t .
0
Each equation contains only four of the five variables and has a different one missing.
It is important to choose the equation that contains the three known variables and one unknown variable for each specific
situation.
Motion diagrams represent the motion of an object by displaying its location at various equally spaced times on the same
diagram.
Motion diagrams show an object’s initial position and velocity and presents several spots in the center of the diagram. These
spots reveal the object’s state of motion.
Motion diagrams contain information about an object’s position at particular time instances and is therefore more informative
than a path diagram.
2.4.2 https://phys.libretexts.org/@go/page/14441
Key Terms
kinematics: The branch of physics concerned with objects in motion.
stroboscopic: Relating to an instrument used to make a cyclically moving object appear to be slow-moving, or stationary.
diagram: A graph or chart.
motion: A change of position with respect to time.
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2.5: Free-Falling Objects
learning objectives
Solve basic problems concerning free fall and distinguish it from other kinds of motion
The motion of falling objects is the simplest and most common example of motion with changing velocity. If a coin and a piece of
paper are simultaneously dropped side by side, the paper takes much longer to hit the ground. However, if you crumple the paper
into a compact ball and drop the items again, it will look like both the coin and the paper hit the floor simultaneously. This is
because the amount of force acting on an object is a function of not only its mass, but also area. Free fall is the motion of a body
where its weight is the only force acting on an object.
s2
. This value is also often expressed as a negative acceleration in mathematical calculations due to the downward direction of
gravity.
The best way to see the basic features of motion involving gravity is to start by considering straight up and down motion with no
air resistance or friction. This means that if the object is dropped, we know the initial velocity is zero. Once the object is in motion,
the object is in free-fall. Under these circumstances, the motion is one-dimensional and has constant acceleration, gg. The
kinematic equations for objects experiencing free fall are:
v = v0 − gt (2.5.1)
1
2
y = y0 + v0 t − gt (2.5.2)
2
2 2
v =v − 2g(y − y0 ), (2.5.3)
0
2.5.1 https://phys.libretexts.org/@go/page/14442
Free Fall Motion
Video 2.5.1 : Free Fall Motion - Describes how to calculate the time for an object to fall if given the height and the height that an
object fell if given the time to fall.
Example 2.5.1:
Key Points
The acceleration of free-falling objects is called the acceleration due to gravity, since objects are pulled towards the center of
the earth.
The acceleration due to gravity is constant on the surface of the Earth and has the value of 9.80 . m
2
s
2.5.2 https://phys.libretexts.org/@go/page/14442
Glossary
Acceleration
The amount by which a speed or velocity changes within a certain period of time (and so a scalar quantity or a vector quantity).
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CHAPTER OVERVIEW
This page titled 20: Circuits and Direct Currents is shared under a not declared license and was authored, remixed, and/or curated by Boundless.
1
20.1: Overview
learning objectives
Describe structure of an electrical circuit and identify elements of a direct current circuit
Overview
An electrical network is an interconnection of electrical elements such as resistors, inductors, capacitors, transmission lines, voltage
sources, current sources and switches. An electrical circuit is a special type of network, one that has a closed loop giving a return
path for the current. Electrical networks that consist only of sources (voltage or current), linear lumped elements (resistors,
capacitors, inductors), and linear distributed elements (transmission lines) can be analyzed by algebraic and transform methods. A
resistive circuit is a circuit containing only resistors and ideal current and voltage sources. Analysis of resistive circuits is less
complicated than analysis of circuits containing capacitors and inductors. If the sources are constant (DC) sources, the result is a
DC circuit. A network that contains active electronic components is known as an electronic circuit. Such networks are generally
nonlinear and require more complex design and analysis tools.
Circuits: A brief introduction to electric circuits and current flow for introductory physics students.
20.1.1 https://phys.libretexts.org/@go/page/14564
Direct Current Circuits
Direct current (DC) is the unidirectional flow of electric charge. Direct current is produced by sources such as batteries,
thermocouples, solar cells, and commutator-type electric machines of the dynamo type. Direct current may flow in a conductor
such as a wire, but can also flow through semiconductors, insulators, or even through a vacuum as in electron or ion beams. The
electric charge flows in a constant direction, distinguishing it from alternating current (AC). A term formerly used for direct current
was galvanic current.
A direct current circuit is an electrical circuit that consists of any combination of constant voltage sources, constant current sources,
and resistors. In this case, the circuit voltages and currents are independent of time. A particular circuit voltage or current does not
depend on the past value of any circuit voltage or current. This implies that the system of equations that represent a DC circuit do
not involve integrals or derivatives with respect to time. If a capacitor or inductor is added to a DC circuit, the resulting circuit is
not, strictly speaking, a DC circuit. However, most such circuits have a DC solution. This solution gives the circuit voltages and
currents when the circuit is in DC steady state. Such a circuit is represented by a system of differential equations. The solution to
these equations usually contain a time varying or transient part as well as constant or steady state part. It is this steady state part that
is the DC solution. In electronics, it is common to refer to a circuit that is powered by a DC voltage source such as a battery or the
output of a DC power supply as a DC circuit even though what is meant is that the circuit is DC powered.
A simple DC circuit is illustrated in. In circuit diagrams such as this, electrical elements are represented by symbols and usually
labeled with appropriate characteristics, such as the resistance r of a resistor. The electric potential and current may also be labeled
at various points of the circuit. However, please be aware that circuit diagram conventions do differ among textbooks and subject
fields, leading to different symbols being used for the same circuit elements.
Example Circuit Element Symbols: A set of example circuit elements and their associated symbols commonly used in circuit
diagrams.
A Simple Circuit: A simple DC circuit with a voltage source (V) and resistor (R). The current i flowing through the circuit is given
by Ohm’s law.
Sources of EMF
Electromotive force (EMF) is the voltage voltage generated by a battery or by the magnetic force according to Faraday’s Law of
Induction.
20.1.2 https://phys.libretexts.org/@go/page/14564
learning objectives
Give examples of the devices that can provide the electromotive force
Electromotive Force
Electromotive force, also called EMF (denoted and measured in volts) refers to voltage generated by a battery or by the magnetic
force according to Faraday’s Law of Induction, which states that a time varying magnetic field will induce an electric current.
Electromotive “force” is not considered a force (as force is measured in newtons) but a potential, or energy per unit of charge,
measured in volts. Formally, EMF is classified as the external work expended per unit of charge to produce an electric potential
difference across two open-circuited terminals. By separating positive and negative charges, electric potential difference is
produced, generating an electric field. The created electrical potential difference drives current flow if a circuit is attached to the
source of EMF. When current flows, however, the voltage across the terminals of the source of EMF is no longer the open-circuit
value, due to voltage drops inside the device due to its internal resistance.
Devices that can provide EMF include electrochemical cells (batteries), thermoelectric devices, solar cells, electrical generators,
transformers, and even Van de Graaff generators (examples shown in ).
Examples of generators of electromotive force.: A variety of voltage sources (clockwise from top left): the Brazos Wind Farm in
Fluvanna, Texas (credit: Leaflet, Wikimedia Commons); the Krasnoyarsk Dam in Russia (credit: Alex Polezhaev); a solar farm
(credit: U.S. Department of Energy); and a group of nickel metal hydride batteries (credit: Tiaa Monto). The voltage output of each
depends on its construction and load, and equals emf only if there is no load.
In the case of a battery, charge separation that gives rise to a voltage difference is accomplished by chemical reactions at the
electrodes; voltaic cells can be thought of as having a “charge pump” of atomic dimensions at each electrode.
In the case of an electrical generator, a time-varying magnetic field inside the generator creates an electric field via electromagnetic
induction, which in turn creates an energy difference between generator terminals. Charge separation takes place within the
generator, with electrons flowing away from one terminal and toward the other, until, in the open-circuit case, sufficient electric
field builds up to make further movement unfavorable. Again the EMF is countered by the electrical voltage due to charge
separation. If a load is attached, this voltage can drive a current. The general principle governing the EMF in such electrical
machines is Faraday’s law of Induction.
In nature, EMF is generated whenever magnetic field fluctuations occur through a surface. An example of this is the varying Earth
magnetic field during a geomagnetic storm, acting on anything on the surface of the planet, like an extended electrical grid.
Key Points
Direct current ( DC ) is the unidirectional flow of electric charge. Direct current is produced by sources such as batteries,
thermocouples, solar cells, and commutator-type electric machines of the dynamo type.
A direct current circuit is an electrical circuit that consists of constant voltage sources, constant current sources, and resistors. It
is common to refer to a circuit that is powered by a DC voltage source as a DC circuit even though what is meant is that the
20.1.3 https://phys.libretexts.org/@go/page/14564
circuit is DC powered.
A number of electrical laws apply to all electrical networks: Ohm ‘s law, Kirchhoff’s current and voltage laws, Norton’s
theorem, and Thévenin’s theorem. With these, it is possible to calculate the electric potential and current at each point in the
circuit.
EMF is classified as the external work expended per unit of charge to produce an electric potential difference across two open-
circuited terminals.
Devices that can provide EMF include electrochemical cells, thermoelectric devices, solar cells, electrical generators,
transformers, and even Van de Graaff generators.
In nature, EMF is generated whenever magnetic field fluctuations occur through a surface.
Key Terms
current-voltage characteristic: A current–voltage characteristic or I–V curve (current–voltage curve) is a relationship between
the electric current through a circuit, device, or material, and the corresponding voltage, or potential difference across it.
electrical circuit: An interconnection of electrical elements such as resistors, inductors, capacitors, transmission lines, voltage
sources, current sources and switches that has a closed loop giving a return path for the current.
direct current: An electric current in which the electrons flow in one direction, but may vary with time.
battery: A device that produces electricity by a chemical reaction between two substances.
electromotive force: (EMF)—The voltage generated by a battery or by the magnetic force according to Faraday’s Law. It is
measured in units of volts, not newtons, and thus, is not actually a force.
Faraday’s law of induction: Is a basic law of electromagnetism that predicts how a magnetic field will interact with an electric
circuit to produce an electromotive force.
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20.2: Resistors in Series and Parallel
learning objectives
Calculate the total resistance in the circuit with resistors connected in series
Overview
Most circuits have more than one component, called a resistor, that limits the flow of charge in the circuit. A measure of this limit
on charge flow is called resistance. The simplest combinations of resistors are the series and parallel connections. The total
resistance of a combination of resistors depends on both their individual values and how they are connected.
Series Circuits: A brief introduction to series circuit and series circuit analysis, including Kirchhoff’s Current Law (KCL) and
Kirchhoff’s Voltage Law (KVL).
Resistors in Series
Resistors are in series whenever the flow of charge, or the current, must flow through components sequentially.
20.2.1 https://phys.libretexts.org/@go/page/14565
Resistors in Series: These four resistors are connected in series because if a current was applied at one end, it would flow through
each resistor sequentially to the end.
shows resistors in series connected to a voltage source. The total resistance in the circuit is equal to the sum of the individual
resistances, since the current has to pass through each resistor in sequence through the circuit.
Resistors connected in a series circuit: Three resistors connected in series to a battery (left) and the equivalent single or series
resistance (right).
equal: V = V + V + V , based on the conservation of energy and charge. If we substitute the values for individual voltages, we
1 2 3
get:
or
V = I (R1 + R2 + R3 ) (20.2.2)
This implies that the total resistance in a series is equal to the sum of the individual resistances. Therefore, for every circuit with N
number of resistors connected in series:
Since all of the current must pass through each resistor, it experiences the resistance of each, and resistances in series simply add
up.
Since voltage and resistance have an inverse relationship, individual resistors in series do not get the total source voltage, but divide
it. This is indicated in an example of when two light bulbs are connected together in a series circuit with a battery. In a simple
circuit consisting of one 1.5V battery and one light bulb, the light bulb would have a voltage drop of 1.5V across it. If two
lightbulbs were connected in series with the same battery, however, they would each have 1.5V/2, or 0.75V drop across them. This
would be evident in the brightness of the lights: each of the two light bulbs connected in series would be half as dim as the single
light bulb. Therefore, resistors connected in series use up the same amount of energy as a single resistor, but that energy is divided
up between the resistors depending on their resistances.
Resistors in Parallel
The total resistance in a parallel circuit is equal to the sum of the inverse of each individual resistances.
20.2.2 https://phys.libretexts.org/@go/page/14565
learning objectives
Calculate the total resistance in the circuit with resistors connected in parallel
Overview
Resistors in a circuit can be connected in series or in parallel. The total resistance of a combination of resistors depends on both
their individual values and how they are connected.
Parallel Circuits: A brief overview of parallel circuit analysis using VIRP tables for high school physics students.
Resistors in Parallel
Resistors are in parallel when each resistor is connected directly to the voltage source by connecting wires having negligible
resistance. Each resistor thus has the full voltage of the source applied to it.
20.2.3 https://phys.libretexts.org/@go/page/14565
Resistors in Parallel: A parallel connection of resistors.
Each resistor draws the same current it would if it were the only resistor connected to the voltage source. This is true of the
circuitry in a house or apartment. Each outlet that is connected to a appliance (the “resistor”) can operate independently, and the
current does not have to pass through each appliance sequentially.
Parallel resistors: Three resistors connected in parallel to a battery and the equivalent single or parallel resistance.
I = I1 + I2 + I3 (20.2.4)
or
1 1 1
I = V( + + ) (20.2.6)
R1 R2 R3
This implies that the total resistance in a parallel circuit is equal to the sum of the inverse of each individual resistances. Therefore,
for every circuit with nn number or resistors connected in parallel,
1 1 1 1
Rn( parallel )
= + + …+ (20.2.7)
R1 R2 R3 Rn
This relationship results in a total resistance that is less than the smallest of the individual resistances. When resistors are connected
in parallel, more current flows from the source than would flow for any of them individually, so the total resistance is lower.
Each resistor in parallel has the same full voltage of the source applied to it, but divide the total current amongst them. This is
exemplified by connecting two light bulbs in a parallel circuit with a 1.5V battery. In a series circuit, the two light bulbs would be
half as dim when connected to a single battery source. However, if the two light bulbs were connected in parallel, they would be
equally as bright as if they were connected individually to the battery. Because the same full voltage is being applied to both light
bulbs, the battery would also die more quickly, since it is essentially supplying full energy to both light bulbs. In a series circuit, the
battery would last just as long as it would with a single light bulb, only the brightness is then divided amongst the bulbs.
20.2.4 https://phys.libretexts.org/@go/page/14565
Combination Circuits
A combination circuit can be broken up into similar parts that are either series or parallel.
learning objectives
Describe arrangement of resistors in a combination circuit and its practical implications
Combination Circuits
More complex connections of resistors are sometimes just combinations of series and parallel. This is commonly encountered,
especially when wire resistances is considered. In that case, wire resistance is in series with other resistances that are in parallel.
A combination circuit can be broken up into similar parts that are either series or parallel, as diagrammed in. In the figure, the total
resistance can be calculated by relating the three resistors to each other as in series or in parallel. R1 and R2 are connected in
parallel in relation to each other, so we know that for that subset, the inverse of resistance would be equal to:
Resistor Network: In this combination circuit, the circuit can be broken up into a series component and a parallel component.
20.2.5 https://phys.libretexts.org/@go/page/14565
Combination Circuits: Two parallel resistors in series with one resistor.
R1 R2
1
R1
+
1
R2
or R1 +R2
R3 is connected in series to both R1 and R2, so the resistance would be calculated as:
R1 R2
R = + R3 (20.2.8)
R1 + R2
Reducing a combination circuit: This combination of seven resistors has both series and parallel parts. Each is identified and
reduced to an equivalent resistance, and these are further reduced until a single equivalent resistance is reached.
Reducing those parallel resistors into a single R value allows us to visualize the circuit in a more simplified manner. In the top right
image, we can see that the circled portion contains two resistors in series. We can further reduce that to another R value by adding
them. The next step shows that the circled two resistors are in parallel. Reducing those highlights that the last two are in series, and
thus can be reduced to a single resistance value for the entire circuit.
One practical implication of a combination circuit is that resistance in wires reduces the current and power delivered to a resistor.
Combination circuit can be transformed into a series circuit, based on an understanding of the equivalent resistance of parallel
branches to a combination circuit. A series circuit can be used to determine the total resistance of the circuit. Essentially, wire
resistance is a series with the resistor. It thus increases the total resistance and decreases the current. If wire resistance is relatively
large, as in a worn (or a very long) extension cord, then this loss can be significant. If a large current is drawn, the IR drop in the
wires can also be significant.
20.2.6 https://phys.libretexts.org/@go/page/14565
Charging a Battery: EMFs in Series and Parallel
When voltage sources are connected in series, their emfs and internal resistances are additive; in parallel, they stay the same.
learning objectives
Compare the resistances and electromotive forces for the voltage sources connected in the same and opposite polarity, and
in series and in parallel
When more than one voltage source is used, they can be connected either in series or in parallel, similar to resistors in a circuit.
When voltage sources are in series facing the same direction, their internal resistances add and their electromotive force, or emf,
add algebraically. These types of voltage sources are common in flashlights, toys, and other appliances. Usually, the cells are in
series in order to produce a larger total emf.
Flashlight and Bulb: A series connection of two voltage sources in the same direction. This schematic represents a flashlight with
two cells (voltage sources) and a single bulb (load resistance) in series.
A battery is a multiple connection of voltaic cells. The disadvantage of series connections of cells in this manner, though, is that
their internal resistances add. This can sometimes be problematic. For example, if you placed two 6v batteries in your car instead of
the typical 12v single battery, you would be adding both the emfs and the internal resistances of each battery. You would therefore
end up with the same 12v emf, though the internal resistance would then be doubled, causing issues for you when you want to start
your engine.
But, if the cells oppose one another—such as when one is put into an appliance backwards—the total emf is less, since it is the
algebraic sum of the individual emfs. When it is reversed, it produces an emf that opposes the other, and results in a difference
between the two voltage sources.
Battery Charger: This represents two voltage sources connected in series with their emfs in opposition. Current flows in the
direction of the greater emf and is limited by the sum of the internal resistances. (Note that each emf is represented by script E in
the figure. ) A battery charger connected to a battery is an example of such a connection. The charger must have a larger emf than
the battery to reverse current through it.
20.2.7 https://phys.libretexts.org/@go/page/14565
When two voltage sources with identical emfs are connected in parallel and also connected to a load resistance, the total emf is the
same as the individual emfs. But the total internal resistance is reduced, since the internal resistances are in parallel. Thus, the
parallel connection can produce a larger current.
Two Identical EMFs: Two voltage sources with identical emfs (each labeled by script E) connected in parallel produce the same
emf but have a smaller total internal resistance than the individual sources. Parallel combinations are often used to deliver more
current.
learning objectives
Express the relationship between the electromotive force and terminal voltage in a form of equation
When you forget to turn off your car lights, they slowly dim as the battery runs down. Why don’t they simply blink off when the
battery’s energy is gone? Their gradual dimming implies that battery output voltage decreases as the battery is depleted. The reason
for the decrease in output voltage for depleted or overloaded batteries is that all voltage sources have two fundamental parts—a
source of electrical energy and an internal resistance.
Electromotive Force
All voltage sources create a potential difference and can supply current if connected to a resistance. On a small scale, the potential
difference creates an electric field that exerts force on charges, causing current. We call this potential difference the electromotive
force (abbreviated emf). Emf is not a force at all; it is a special type of potential difference of a source when no current is flowing.
Units of emf are volts.
Electromotive force is directly related to the source of potential difference, such as the particular combination of chemicals in a
battery. However, emf differs from the voltage output of the device when current flows. The voltage across the terminals of a
battery, for example, is less than the emf when the battery supplies current, and it declines further as the battery is depleted or
loaded down. However, if the device’s output voltage can be measured without drawing current, then output voltage will equal emf
(even for a very depleted battery).
20.2.8 https://phys.libretexts.org/@go/page/14565
Terminal Voltage
presents a schematic representation of a voltage source. The voltage output of a device is measured across its terminals and is
called its terminal voltage V. Terminal voltage is given by the equation:
Schematic Representation of a Voltage Source: Any voltage source (in this case, a carbon-zinc dry cell) has an emf related to its
source of potential difference, and an internal resistance r related to its construction. (Note that the script E stands for emf. ) Also
shown are the output terminals across which the terminal voltage V is measured. Since V=emf−Ir, terminal voltage equals emf only
if there is no current flowing.
where r is the internal resistance and I is the current flowing at the time of the measurement.
I is positive if current flows away from the positive terminal. The larger the current, the smaller the terminal voltage. Likewise, it is
true that the larger the internal resistance, the smaller the terminal voltage.
Key Points
The same current flows through each resistor in series.
Individual resistors in series do not get the total source voltage, but divide it.
The total resistance in a series circuit is equal to the sum of the individual resistances:
RN( series ) = R + R + R + … + R
1 2 3 N .
The total resistance in a parallel circuit is less than the smallest of the individual resistances.
Each resistor in parallel has the same voltage of the source applied to it (voltage is constant in a parallel circuit).
Parallel resistors do not each get the total current; they divide it (current is dependent on the value of each resistor and the
number of total resistors in a circuit).
More complex connections of resistors are sometimes just combinations of series and parallel.
Various parts of a combination circuit can be identified as series or parallel, reduced to their equivalents, and then further
reduced until a single resistance is left.
Resistance in wires reduces the current and power delivered to a resistor. If the resistance in wires is relatively large, as in a
worn (or a very long) extension cord, then this loss can be significant and affect power output into appliances.
Emfs connected in the same polarity in series are additive and result in a higher total emf.
Two emfs connected in the opposite polarity in series have a total emf equal to the difference between them, and can be used to
charge the lower voltage source.
20.2.9 https://phys.libretexts.org/@go/page/14565
Two voltage sources with identical emfs connected in parallel have a net emf equivalent to one emf source, however, the net
internal resistance is less, and therefore produces a higher current.
Key Terms
series: A number of things that follow on one after the other or are connected one after the other.
resistance: The opposition to the passage of an electric current through that element.
parallel: An arrangement of electrical components such that a current flows along two or more paths.
combination circuit: An electrical circuit containing multiple resistors that are connected in a combination of both series and
parallel connections.
electromotive force: (EMF)—The voltage generated by a battery or by the magnetic force according to Faraday’s Law. It is
measured in units of volts, not newtons, and thus, is not actually a force.
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20.3: Kirchhoff’s Rules
learning objectives
Describe relationship between the Kirchhoff’s circuit laws and the energy and charge in the electrical circuits
Closed Circuit: To determine all variables (i.e., current and voltage drops across the different resistors) in this circuit, Kirchhoff’s
rules must be applied.
As a final note, Kirchhoff’s laws depend on certain conditions. The voltage law is a simplification of Faraday’s law of induction,
and is based on the assumption that there is no fluctuating magnetic field within the closed loop. Thus, although this law can be
applied to circuits containing resistors and capacitors (as well as other circuit elements), it can only be used as an approximation to
the behavior of the circuit when a changing current and therefore magnetic field are involved.
learning objectives
Describe relationship between the Kirchhoff’s circuit laws and the energy and charge in the electrical circuits
Kirchhoff’s junction rule, also known as Kirchhoff’s current law (KCL), Kirchoff’s first law, Kirchhoff’s point rule, and
Kirchhoff’s nodal rule, is an application of the principle of conservation of electric charge.
Kirchhoff’s junction rule states that at any junction ( node ) in an electrical circuit, the sum of the currents flowing into that
junction is equal to the sum of the currents flowing out of that junction. In other words, given that a current will be positive or
negative depending on whether it is flowing towards or away from a junction, the algebraic sum of currents in a network of
conductors meeting at a point is equal to zero. A visual representation can be seen in.
20.3.1 https://phys.libretexts.org/@go/page/14566
Kirchhoff’s Junction Law: Kirchhoff’s Junction Law illustrated as currents flowing into and out of a junction.
Kirchhoff’s Loop and Junction Rules Theory: We justify Kirchhoff’s Rules from diarrhea and conservation of energy. Some
people call ’em laws, but not me!
Thus, Kirchoff’s junction rule can be stated mathematically as a sum of currents (I):
n
∑ Ik = 0 (20.3.1)
k=1
where n is the total number of branches carrying current towards or away from the node.
20.3.2 https://phys.libretexts.org/@go/page/14566
This law is founded on the conservation of charge (measured in coulombs), which is the product of current (amperes) and time
(seconds).
Limitation
Kirchhoff’s junction law is limited in its applicability. It holds for all cases in which total electric charge (Q) is constant in the
region in consideration. Practically, this is always true so long as the law is applied for a specific point. Over a region, however,
charge density may not be constant. Because charge is conserved, the only way this is possible is if there is a flow of charge across
the boundary of the region. This flow would be a current, thus violating Kirchhoff’s junction law.
learning objectives
Formulate the Kirchhoff’s loop rule, noting its assumptions
Kirchhoff’s loop rule (otherwise known as Kirchhoff’s voltage law (KVL), Kirchhoff’s mesh rule, Kirchhoff’s second law,
orKirchhoff’s second rule) is a rule pertaining to circuits, and is based on the principle of conservation of energy.
Conservation of energy—the principle that energy is neither created nor destroyed—is a ubiquitous principle across many studies
in physics, including circuits. Applied to circuitry, it is implicit that the directed sum of the electrical potential differences
(voltages) around any closed network is equal to zero. In other words, the sum of the electromotive force (emf) values in any closed
loop is equal to the sum of the potential drops in that loop (which may come from resistors).
Another equivalent statement is that the algebraic sum of the products of resistances of conductors (and currents in them) in a
closed loop is equal to the total electromotive force available in that loop. Mathematically, Kirchhoff’s loop rule can be represented
as the sum of voltages in a circuit, which is equated with zero:
20.3.3 https://phys.libretexts.org/@go/page/14566
Kirchhoff’s Loop and Junction Rules Theory: We justify Kirchhoff’s Rules from diarrhea and conservation of energy. Some
people call ’em laws, but not me!
n
∑ Vk = 0 (20.3.2)
k=1
Here, Vk is the voltage across element k, and n is the total number of elements in the closed loop circuit. An illustration of such a
circuit is shown in. In this example, the sum of v1, v2, v3, and v4 (and v5 if it is included), is zero.
Kirchhoff’s Loop Rule: Kirchhoff’s loop rule states that the sum of all the voltages around the loop is equal to zero: v1 + v2 + v3
– v4 = 0.
Given that voltage is a measurement of energy per unit charge, Kirchhoff’s loop rule is based on the law of conservation of energy,
which states: the total energy gained per unit charge must equal the amount of energy lost per unit of charge.
Example 20.3.1:
illustrates the changes in potential in a simple series circuit loop. Kirchhoff’s second rule requires emf−Ir−IR1−IR2=0.
Rearranged, this is emf=Ir+IR1+IR2, meaning that the emf equals the sum of the IR (voltage) drops in the loop. The emf
supplies 18 V, which is reduced to zero by the resistances, with 1 V across the internal resistance, and 12 V and 5 V across the
two load resistances, for a total of 18 V.
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The Loop Rule: An example of Kirchhoff’s second rule where the sum of the changes in potential around a closed loop must
be zero. (a) In this standard schematic of a simple series circuit, the emf supplies 18 V, which is reduced to zero by the
resistances, with 1 V across the internal resistance, and 12 V and 5 V across the two load resistances, for a total of 18 V. (b)
This perspective view represents the potential as something like a roller coaster, where charge is raised in potential by the emf
and lowered by the resistances. (Note that the script E stands for emf.)
Limitation
Kirchhoff’s loop rule is a simplification of Faraday’s law of induction, and holds under the assumption that there is no fluctuating
magnetic field linking the closed loop. In the presence of a variable magnetic field, electric fields could be induced and emf could
be produced, in which case Kirchhoff’s loop rule breaks down.
Applications
Kirchhoff’s rules can be used to analyze any circuit and modified for those with EMFs, resistors, capacitors and more.
learning objectives
Describe conditions when the Kirchoff’s rules are useful to apply
Overview
Kirchhoff’s rules can be used to analyze any circuit by modifying them for those circuits with electromotive forces, resistors,
capacitors and more. Practically speaking, however, the rules are only useful for characterizing those circuits that cannot be
simplified by combining elements in series and parallel.
20.3.5 https://phys.libretexts.org/@go/page/14566
Combinations in series and parallel are typically much easier to perform than applying either of Kirchhoff’s rules, but Kirchhoff’s
rules are more broadly applicable and should be used to solve problems involving complex circuits that cannot be simplified by
combining circuit elements in series or parallel.
Kirchhoff’s Rules: sample problem: This image shows a very complicated circuit, which can be reduced and solved using
Kirchoff’s Rules.
Applying Kirchhoff’s junction rule at point a, we find:
I1 = I2 + I3 (20.3.3)
because I1 flows into point a, while I2 and I3 flow out. The same can be found at point e. We now must solve this equation for each
of the three unknown variables, which will require three different equations.
Considering loop abcdea, we can use Kirchhoff’s loop rule:
Substituting values of resistance and emf from the figure diagram and canceling the ampere unit gives:
−3 I2 + 18 − 6 I1 = 0 (20.3.5)
This is the second part of a system of three equations that we can use to find all three current values. The last can be found by
applying the loop rule to loop aefgha, which gives:
In this case, the signs were reversed compared with the other loop, because elements are traversed in the opposite direction.
We now have three equations that can be used in a system. The second will be used to define I2, and can be rearranged to:
I2 = 6 − 2 I1 (20.3.8)
The third equation can be used to define I3, and can be rearranged to:
I3 = 22.5 − 3 I1 (20.3.9)
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Substituting the new definitions of I2 and I3 (which are both in common terms of I1), into the first equation (I1=I2+I3), we get:
I I1 = (6 − 2 I1 ) + (22.5 − 3 I1 ) = 28.5 − 5 I1 (20.3.10)
Simplifying, we find that I1=4.75 A. Inserting this value into the other two equations, we find that I2=-3.50 A and I3=8.25 A.
Key Points
Kirchhoff used Georg Ohm ‘s work as a foundation to create Kirchhoff’s current law (KCL) and Kirchhoff’s voltage law (KVL)
in 1845. These can be derived from Maxwell’s Equations, which came 16-17 years later.
It is impossible to analyze some closed-loop circuits by simplifying as a sum and/or series of components. In these cases,
Kirchhoff’s laws can be used.
Kirchhoff’s laws are special cases of conservation of energy and charge.
Kirchhoff’s junction rule is an application of the principle of conservation of electric charge: current is flow of charge per time,
and if current is constant, that which flows into a point in a circuit must equal that which flows out of it.
The mathematical representation of Kirchhoff’s law is: ∑ n
k=1
I = 0 where Ik is the current of k, and n is the total number of
k
Kirchhoff’s loop rule is a simplification of Faraday’s law of induction and holds under the assumption that there is no
fluctuating magnetic field linking the closed loop.
Kirchhoff’s rules can be applied to any circuit, regardless of its composition and structure.
Because combining elements is often easy in parallel and series, it is not always convenient to apply Kirchhoff’s rules.
To solve for current in a circuit, the loop and junction rules can be applied. Once all currents are related by the junction rule,
one can use the loop rule to obtain several equations to use as a system to find each current value in terms of other currents.
These can be solved as a system.
Key Terms
resistor: An electric component that transmits current in direct proportion to the voltage across it.
electromotive force: (EMF)—The voltage generated by a battery or by the magnetic force according to Faraday’s Law. It is
measured in units of volts (not newtons, N; EMF is not a force).
capacitor: An electronic component consisting of two conductor plates separated by empty space (sometimes a dielectric
material is instead sandwiched between the plates), and capable of storing a certain amount of charge.
electric charge: A quantum number that determines the electromagnetic interactions of some subatomic particles; by
convention, the electron has an electric charge of -1 and the proton +1, and quarks have fractional charge.
current: The time rate of flow of electric charge.
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20.4: Voltmeters and Ammeters
learning objectives
Compare circuit connection of an ammeter and a voltmeter
Voltmeters and ammeters measure the voltage and current, respectively, of a circuit. Some meters in automobile dashboards, digital
cameras, cell phones, and tuner-amplifiers are voltmeters or ammeters.
Voltmeters and Ammeters: A brief introduction to voltmeters and ammeters for introductory physics students.
Voltmeters
A voltmeter is an instrument that measures the difference in electrical potential between two points in an electric circuit. An analog
voltmeter moves a pointer across a scale in proportion to the circuit’s voltage; a digital voltmeter provides a numerical display. Any
measurement that can be converted to voltage can be displayed on a meter that is properly calibrated; such measurements include
pressure, temperature, and flow.
20.4.1 https://phys.libretexts.org/@go/page/14567
Voltmeter: Demonstration voltmeter from a physics class
In order for a voltmeter to measure a device’s voltage, it must be connected in parallel to that device. This is necessary because
objects in parallel experience the same potential difference.
Voltmeter in Parallel: (a) To measure the potential difference in this series circuit, the voltmeter (V) is placed in parallel with the
voltage source or either of the resistors. Note that terminal voltage is measured between points a and b. It is not possible to connect
the voltmeter directly across the EMF without including its internal resistance, r. (b) A digital voltmeter in use
Ammeters
An ammeter measures the electric current in a circuit. The name is derived from the name for the SI unit for electric current,
amperes (A).
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In order for an ammeter to measure a device’s current, it must be connected in series to that device. This is necessary because
objects in series experience the same current. They must not be connected to a voltage source — ammeters are designed to work
under a minimal burden, (which refers to the voltage drop across the ammeter, typically a small fraction of a volt).
Ammeter in Series: An ammeter (A) is placed in series to measure current. All of the current in this circuit flows through the
meter. The ammeter would have the same reading if located between points d and e or between points f and a, as it does in the
position shown. (Note that the script capital E stands for EMF, and r stands for the internal resistance of the source of potential
difference. )
Galvanometers as Voltmeters
A galvanometer can function as a voltmeter when it is connected in series with a large resistance R. The value of R is determined
by the maximum voltage that will be measured. Suppose you want 10 V to produce a full-scale deflection of a voltmeter containing
a 25-Ω galvanometer with a 50-μA sensitivity. Then 10 V applied to the meter must produce a current of 50 μA. The total
resistance must be:
V 10V
Rtot = R + r = = = 200kΩ (20.4.1)
I 50μA
or:
(R is so large that the galvanometer resistance, r, is nearly negligible. ) Note that 5 V applied to this voltmeter produces a half-scale
deflection by sending a 25-μA current through the meter, and so the voltmeter’s reading is proportional to voltage, as desired. This
voltmeter would not be useful for voltages less than about half a volt, because the meter deflection would be too small to read
accurately. For other voltage ranges, other resistances are placed in series with the galvanometer. Many meters allow a choice of
scales, which involves switching an appropriate resistance into series with the galvanometer.
20.4.3 https://phys.libretexts.org/@go/page/14567
Galvanometers as Ammeters
The same galvanometer can also function as an ammeter when it is placed in parallel with a small resistance R, often called the
shunt resistance. Since the shunt resistance is small, most of the current passes through it, allowing an ammeter to measure currents
much greater than those that would produce a full-scale deflection of the galvanometer.
Suppose, for example, we need an ammeter that gives a full-scale deflection for 1.0 A and that contains the same 25-Ω
galvanometer with 50-μA sensitivity. Since R and r are in parallel, the voltage across them is the same.
These IR drops are: IR = IGr
IG
so that: IR = I
=
R
r
.
Solving for R, and noting that IG is 50 μA and I is 0.999950 A, we have:
IG 50μA −3
R =r = (25Ω) = 1.25 × 10 Ω (20.4.3)
I 0.999950A
Null Measurements
Null measurements balance voltages so there is no current flowing through the measuring devices that would interfere with the
measurement.
learning objectives
Explain why null measurements are employed
Null Measurements
Standard measurements of voltage and current alter circuits, introducing numerical uncertainties. Voltmeters draw some extra
current, whereas ammeters reduce current flow. Null measurements balance voltages, so there is no current flowing through the
measuring device and the circuit is unaltered. Null measurements are generally more accurate but more complex than standard
voltmeters and ammeters. Their precision is still limited.
The Potentiometer
When measuring the EMF of a battery and connecting the battery directly to a standard voltmeter, as shown in, the actual quantity
measured is the terminal voltage V. Voltage is related to the EMF of the battery by V=emf−Ir, where I is the current that flows and r
is the internal resistance of the battery.
Voltmeter Connected to Battery: An analog voltmeter attached to a battery draws a small but nonzero current and measures a
terminal voltage that differs from the EMF of the battery. (Note that the script capital E symbolizes electromotive force, or EMF. )
Since the internal resistance of the battery is not known precisely, it is not possible to calculate the EMF precisely.
The EMF could be accurately calculated if r were known, which is rare. If the current I could be made zero, then V=emf, and EMF
could be directly measured. However, standard voltmeters need a current to operate.
A potentiometer is a null measurement device for measuring potentials (voltages). A voltage source is connected to resistor R,
passing a constant current through it. There is a steady drop in potential (IR drop) along the wire, so a variable potential is obtained
through contact along the wire.
20.4.4 https://phys.libretexts.org/@go/page/14567
An unknown emfx (represented by script Ex) connected in series with a galvanometer is shown in. Note that emfx opposes the other
voltage source. The location of the contact point is adjusted until the galvanometer reads zero. When the galvanometer reads zero,
emfx=IRx, where Rxis the resistance of the wire section up to the contact point. Since no current flows through the galvanometer,
none flows through the unknown EMF, and emfx is sensed.
Potentiometer: The potentiometer is a null measurement device. (a. ) A voltage source connected to a long wire resistor passes a
constant current I through it. (b.) An unknown EMF (labeled script Ex) is connected as shown, and the point of contact along R is
adjusted until the galvanometer reads zero. The segment of wire has a resistance Rx and script Ex=IRx, where I is unaffected by the
connection, since no current flows through the galvanometer. The unknown EMF is thus proportional to the resistance of the wire
segment.
Standard EMF is substituted for emfx, and the contact point is adjusted until the galvanometer reads zero, so that emfs=IRs. In both
cases, no current passes through the galvanometer. The current Ithrough the long wire is identical. Taking the ratio emfx/emfs,
Icancels, and solving for emfx gives what is seen in.
Because a long uniform wire is used for R, the ratio of resistances Rx/Rs is the same as the ratio of the lengths of wire that zero the
galvanometer for each EMF. The three quantities on the right-hand side of the equation are now known or measured, and emfx can
be calculated. There is often less uncertainty in this calculation than when using a voltmeter directly, but it is not zero. There is
always some uncertainty in the ratio of resistances Rx/Rsand in the standard EMFs. Furthermore, it is not possible to tell when the
galvanometer reads exactly zero, which introduces error into both Rx and Rs, and may also affect the current I.
Resistance Measurements
Many so-called ohmmeters measure resistance. Most common ohmmeters apply a voltage to a resistance, measure the current, and
calculate the resistance using Ohm ‘s law. Their readout is this calculated resistance. Simple configurations using standard
voltmeters and ammeters have limited accuracy, because the meters alter both the voltage applied to the resistor and the current
flowing through it. The Wheatstone bridge is a null measurement device for calculating resistance by balancing potential drops in a
circuit. The device is called a bridge because the galvanometer forms a bridge between two branches. A variety of bridge
devicesare used to make null measurements in circuits. Resistors R1 and R2 are precisely known, while the arrow through
R3indicates that it is a variable resistance. The value of R3 can be precisely read. With the unknown resistance Rx in the circuit, R3
is adjusted until the galvanometer reads zero.
20.4.5 https://phys.libretexts.org/@go/page/14567
Wheatstone Bridge: The Wheatstone bridge is used to calculate unknown resistances. The variable resistance R3 is adjusted until
the galvanometer reads zero with the switch closed. This simplifies the circuit, allowing Rx to be calculated based on the IR drops.
The potential difference between points b and d is then zero, meaning that b and d are at the same potential. With no current
running through the galvanometer, it has no effect on the rest of the circuit. So the branches abc and adc are in parallel, and each
branch has the full voltage of the source. Since b and d are at the same potential, the IR drop along ad must equal the IR drop along
ab. Again, since b and d are at the same potential, the IR drop along dcmust equal the IR drop along bc. This equation is used to
calculate the unknown resistance when current through the galvanometer is zero. This method can be very accurate, but it is limited
by two factors. First, it is not possible for the current through the galvanometer to be exactly zero. Second, there are always
uncertainties in R1, R2, and R3, which contribute to the uncertainty in Rx.
Key Points
A voltmeter is an instrument used for measuring electrical potential difference between two points in an electric circuit.
An ammeter is a measuring device used to measure the electric current in a circuit.
A voltmeter is connected in parallel with a device to measure its voltage, while an ammeter is connected in series with a device
to measure its current.
At the heart of most analog meters is a galvanometer, an instrument that measures current flow using the movement, or
deflection, of a needle. The needle deflection is produced by a magnetic force acting on a current-carrying wire.
Measurements of voltages and current with standard voltmeters and ammeters alter the circuit being measured, introducing
uncertainties. Voltmeters draw some extra current, whereas ammeters reduce current flow.
Null measurements are employed to reduce the uncertainty in the measured voltage and current.
The potentiometer and the Wheatstone bridge are two methods for making null measurements.
The potentiometer is an instrument that measures an unknown voltage by opposing with a known voltage, without drawing
current from the voltage source being measured.
A Wheatstone bridge is an electrical circuit used to measure an unknown electrical resistance by balancing two legs of a bridge
circuit, one leg of which includes the unknown component.
Key Terms
shunt resistance: a small resistance R placed in parallel with a galvanometer G to produce an ammeter; the larger the current to
be measured, the smaller R must be; most of the current flowing through the meter is shunted through R to protect the
galvanometer
galvanometer: An analog measuring device, denoted by G, that measures current flow using a needle deflection caused by a
magnetic field force acting upon a current-carrying wire.
null measurements: methods of measuring current and voltage more accurately by balancing the circuit so that no current
flows through the measurement device
20.4.6 https://phys.libretexts.org/@go/page/14567
potentiometer: an instrument that measures a voltage by opposing it with a precise fraction of a known voltage, and without
drawing current from the unknown source.
Wheatstone bridge: An instrument used to measure an unknown electrical resistance by balancing two legs of a bridge circuit,
one leg of which includes the unknown component.
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20.5: RC Circuits
learning objectives
Describe the components and function of an RC circut, noting especially the time-dependence of the capacitor’s charge
An RC circuit is one containing a resistor R and a capacitor C. The capacitor is an electrical component that houses electric charge.
In this Atom, we will study how a series RC circuit behaves when connected to a DC voltage source. (In subsequent Atoms, we
will study its AC behavior. )
Charging
Fig 1 shows a simple RC circuit that employs a DC voltage source. The capacitor is initially uncharged. As soon as the switch is
closed, current flows to and from the initially uncharged capacitor. As charge increases on the capacitor plates, there is increasing
opposition to the flow of charge by the repulsion of like charges on each plate.
Charging an RC Circuit: (a) An RC circuit with an initially uncharged capacitor. Current flows in the direction shown as soon as
the switch is closed. Mutual repulsion of like charges in the capacitor progressively slows the flow as the capacitor is charged,
stopping the current when the capacitor is fully charged and Q=C⋅emf. (b) A graph of voltage across the capacitor versus time, with
the switch closing at time t=0. (Note that in the two parts of the figure, the capital script E stands for emf, q stands for the charge
stored on the capacitor, and τ is the RC time constant. )
In terms of voltage, across the capacitor voltage is given by Vc=Q/C, where Q is the amount of charge stored on each plate and C is
the capacitance. This voltage opposes the battery, growing from zero to the maximum emf when fully charged. Thus, the current
decreases from its initial value of I0=emf/R to zero as the voltage on the capacitor reaches the same value as the emf. When there is
no current, there is no IR drop, so the voltage on the capacitor must then equal the emf of the voltage source.
Initially, voltage on the capacitor is zero and rises rapidly at first since the initial current is a maximum. Fig 1 (b) shows a graph of
capacitor voltage versus time (t) starting when the switch is closed at t=0. The voltage approaches emf asymptotically since the
closer it gets to emf the less current flows. The equation for voltage versus time when charging a capacitor C through a resistor R,
is:
t/RC
V(t) = emf(1 − e ) (20.5.1)
where V(t) is the voltage across the capacitor and emf is equal to the emf of the DC voltage source. (The exact form can be derived
by solving a linear differential equation describing the RC circuit, but this is slightly beyond the scope of this Atom. ) Note that the
unit of RC is second. We define the time constant τ for an RC circuit as τ = RC . τ shows how quickly the circuit charges or
discharges.
Discharging
Discharging a capacitor through a resistor proceeds in a similar fashion, as illustrates. Initially, the current is I0=V0/R, driven by the
initial voltage V0 on the capacitor. As the voltage decreases, the current and hence the rate of discharge decreases, implying another
exponential formula for V. Using calculus, the voltage V on a capacitor C being discharged through a resistor R is found to be
−t/RC
V(t) = V0 e (20.5.2)
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Impedance
Impedance is the measure of the opposition that a circuit presents to the passage of a current when a voltage is applied.
learning objectives
Express the relationship between the impedance, the resistance, and the capacitance of a series RC circuit in a form of
equation
Rather than solving the differential equation relating to circuits that contain resistors and capacitors, we can imagine all sources in
the circuit are complex exponentials having the same frequency. This technique is useful in solving problems in which phase
relationship is important. The phase of the complex impedance is the phase shift by which the current is ahead of the voltage.
Complex Analysis
For an RC circuit in, the AC source driving the circuit is given as:
where V is the amplitude of the AC voltage, j is the imaginary unit (j2=-1), and ωω is the angular frequency of the AC source. Two
things to note:
1. We use lower case alphabets for voltages and sources to represent that they are alternating (i.e., we use vin(t) instead of Vin(t)).
2. The imaginary unit is given the symbol “j”, not the usual “i”. “i” is reserved for alternating currents.
Complex Impedance
The advantage of assuming sources take this form is that all voltages and currents in the circuit are also complex exponentials
(having the same frequency as the source). To appreciate the reason for this, we can investigate how each circuit element behaves
when either the voltage or current is a complex exponential. For the resistor, v = Ri . From our voltage given above, i = e . V
R
jωt
Thus the resistor’s voltage is a complex, as is the current with an amplitude I = . For a capacitor, i = C . Letting the voltage
V
R
dv
dt
be a complex exponential we have i = jωCVe . The amplitude of this complex exponential is I = jωCV .
jωt
The major consequence of assuming complex exponential voltage and currents is that the ratio Z = for rather than depending on
V
time each element depends on source frequency. This quantity is known as the element’s (complex) impedance. The magnitude of
the complex impedance is the ratio of the voltage amplitude to the current amplitude. Just like resistance in DC cases, impedance is
the measure of the opposition that a circuit presents to the passage of a current when a voltage is applied. The impedance of a
resistor is R, while that of a capacitor (C) is . In the case of the circuit in, to find the complex impedance of the RC circuit, we
1
jωC
add the impedance of the two components, just as with two resistors in series: Z = R + 1
jωC
.
(real value) impedance is the real part of the complex impedance Z. For a series RC circuit, we get Z = √R
2
+(
1
ωC
) . We see
that the amplitude of the current will be V
Z
=
V
2
.
2 1
√R +( )
ωC
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learning objectives
Compare the currents in the resistor and capacitor in a series RC circuit connected to an AC voltage source
Phase Angle
Impedance is an AC (alternating current) analogue to resistance in a DC circuit. As we studied in a previously Atom
(“Impedance”), current, voltage and impedance in an RC circuit are related by an AC version of Ohm ‘s law: I = , where I and V V
are peak current and peak voltage respectively, and Z is the impedance of the circuit.
In a series RC circuit connected to an AC voltage source as shown in, conservation of charge requires current be the same in each
part of the circuit at all times. Therefore we can say: the currents in the resistor and capacitor are equal and in phase. (We will
represent instantaneous current as i(t). )
1
= i(t)R + i(t)
jωC
1
= i(t) (R + )
jωC
where ωω is the angular frequency of the AC voltage source and j is the imaginary unit; j2=-1. Since the complex number
−−−−−−−−−−
2
Z =R+
1
jωC
= √R
2
+(
1
ωC
)
jϕ
e has a phase angle ϕ that satisfies cos ϕ = R
2
,
2 1
√R +( )
ωC
we notice that voltage v(t) and current i(t) has a phase difference of ϕ .
For R = 0, ϕ = 90 . As learned from the preceding series of Atoms—the voltage across the capacitor VC follows the current by
∘
Power Factor
Because voltage and current are out of phase, power dissipated by the circuit is not equal to: (peak voltage) times (peak current).
The fact that source voltage and current are out of phase affects the power delivered to the circuit. It can be shown that the average
power is IrmsVrmscosϕ, where Irms and Vrms are the root mean square (rms) averages of the current and voltage, respectively. For
this reason, cosϕ is called the power factor, which can range from 0 to 1.
Key Points
In an RC circuit connected to a DC voltage source, the current decreases from its initial value of I0=emf/R to zero as the voltage
on the capacitor reaches the same value as the emf.
In an RC circuit connected to a DC voltage source, voltage on the capacitor is initially zero and rises rapidly at first since the
initial current is a maximum: V(t) = emf(1 − e ) .
t/RC
The time constant τ for an RC circuit is defined to be RC. It’s unit is in seconds and shows how quickly the circuit charges or
discharges.
The advantage of assuming that sources have complex exponential form is that all voltages and currents in the circuit are also
complex exponentials, having the same frequency as the source.
The major consequence of assuming complex exponential voltage and currents is that the ratio (Z = V/I) for each element does
not depend on time, but does depend on source frequency.
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−−−−−−−−−−
2
ωC
) .
In a series RC circuit connected to an AC voltage source, the currents in the resistor and capacitor are equal and in phase.
In a series RC circuit connected to an AC voltage source, the total voltage should be equal to the sum of voltages on the resistor
and capacitor.
In a series RC circuit connected to an AC voltage source, voltage and current have a phase difference of ϕ , where
cos ϕ =
R
. cosϕ is called the power factor.
2
2 1
√R +( )
ωC
Key Terms
DC: Direct current; the unidirectional flow of electric charge.
capacitor: An electronic component capable of storing an electric charge, especially one consisting of two conductors separated
by a dielectric.
differential equation: An equation involving the derivatives of a function.
impedance: A measure of the opposition to the flow of an alternating current in a circuit; the aggregation of its resistance,
inductive and capacitive reactance. Represented by the symbol Z.
resistor: An electric component that transmits current in direct proportion to the voltage across it.
alternating current: (AC)—An electric current in which the direction of flow of the electrons reverses periodically having an
average of zero, with positive and negative values (with a frequency of 50 Hz in Europe, 60 Hz in the US, 400 Hz for airport
lighting, and some others); especially such a current produced by a rotating generator or alternator.
rms: Root mean square: a statistical measure of the magnitude of a varying quantity.
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CHAPTER OVERVIEW
21: Magnetism
Topic hierarchy
21.1: Magnetism and Magnetic Fields
21.2: Magnets
21.3: Magnetic Force on a Moving Electric Charge
21.4: Motion of a Charged Particle in a Magnetic Field
21.5: Magnetic Fields, Magnetic Forces, and Conductors
21.6: Applications of Magnetism
21: Magnetism is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
21.1: Magnetism and Magnetic Fields
learning objectives
Describe shape of a magnetic field produced by an electric current flowing through a wire
Magnetic Field Generated by Current: (a) Compasses placed near a long straight current-carrying wire indicate that field lines
form circular loops centered on the wire. (b) Right hand rule 2 states that, if the right hand thumb points in the direction of the
current, the fingers curl in the direction of the field. This rule is consistent with the field mapped for the long straight wire and is
valid for any current segment.
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High School Physics - Magnets and Magnetic Fields
Magnets and Magnetic Fields: A brief introduction to magnetism for introductory physics students.
For a long straight wire where I is the current, r is the shortest distance to the wire, and the constant 0=4π10−7 T⋅m/A is the
permeability of free space. (μ0 is one of the basic constants in nature, related to the speed of light. ) Since the wire is very long, the
magnitude of the field depends only on distance from the wire r, not on position along the wire. This is one of the simplest cases to
calculate the magnetic field strenght from a current.
The magnetic field of a long straight wire has more implications than one might first suspect. Each segment of current produces a
magnetic field like that of a long straight wire, and the total field of any shape current is the vector sum of the fields due to each
segment. The formal statement of the direction and magnitude of the field due to each segment is called the Biot-Savart law.
Integral calculus is needed to sum the field for an arbitrary shape current. The Biot-Savart law is written in its complete form as:
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μ0 I dl × ^
r
B= ∫ (21.1.2)
4π r2
where the integral sums over the wire length where vector dℓ is the direction of the current; r is the distance between the location of
dℓ, and the location at which the magnetic field is being calculated; and r ̂ is a unit vector in the direction of r. The reader may apply
the simplifications in calculating the magnetic field from an infinite straight wire as above and see that the Biot-Savart law reduces
to the first, simpler equation.
Ampere’s Law
A more fundamental law than the Biot-Savart law is Ampere’s Law, which relates magnetic field and current in a general way. In
SI units, the integral form of the original Ampere’s circuital law is a line integral of the magnetic field around some closed curve C
(arbitrary but must be closed). The curve C in turn bounds both a surface S through which the electric current passes through (again
arbitrary but not closed—since no three-dimensional volume is enclosed by S), and encloses the current. You can think of the
“surface” as the cross-sectional area of a wire carrying current.
The mathematical statement of the law states that the total magnetic field around some path is directly proportional to the current
which passes through that enclosed path. It can be written in a number of forms, one of which is given below.
∮ B ⋅ dl = μ0 ∬ J ⋅ dS = μ0 Ienc (21.1.3)
S
where the magnetic field is integrated over a curve (circumfrence of a wire), equivalent to integrating the current density (in
amperes per square meter, Am-2) over the cross section area of the wire (which is equal to the permeability constant times the
enclosed current Ienc). Ampere’s law is always valid for steady currents and can be used to calculate the B-field for certain highly
symmetric situations such as an infinite wire or an infinite solenoid. Ampere’s Law is also a component of Maxwell’s Equations.
Force on a Current-Carrying Wire: The right hand rule can be used to determine the direction of the force on a current-carrying
wire placed in an external magnetic field.
F = qvB sin θ (21.1.4)
where i = nqvA . The right hand rule can give you the direction of the force on the wire, as seen in the above figure. Note that the
B-field in this case is the external field.
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Permanent Magnets
Permanent magnets are objects made from ferromagnetic material that produce a persistent magnetic field.
learning objectives
Give examples and counterexamples of permanent magnets
Permanent Magnets
Overview
Recall that a magnet is a material or object that generates a magnetic field. This magnetic field is invisible but is responsible for the
most notable property of a magnet: a force that pulls on other ferromagnetic materials, such as iron, and attracts or repels other
magnets.
Types of Magnets
A permanent magnet is an object made from a material that is magnetized and creates its own persistent magnetic field. An
everyday example is a refrigerator magnet used to hold notes on a refrigerator door. Materials that can be magnetized, which are
also the ones that are strongly attracted to a magnet, are called ferromagnetic. These include iron, nickel, cobalt, some alloys of rare
earth metals, and some naturally occurring minerals such as lodestone. Although ferromagnetic materials are the only ones
attracted to a magnet strongly enough to be commonly considered magnetic, all other substances respond weakly to a magnetic
field, by one of several other types of magnetism. The counterexample to a permanent magnet is an electromagnet, which only
becomes magnetized when an electric current flows through it.
Example of a Permanent Magnet: An example of a permanent magnet: a “horseshoe magnet” made of alnico, an iron alloy. The
magnet is made in the shape of a horseshoe to bring the two magnetic poles close to each other, to create a strong magnetic field
there that can pick up heavy pieces of iron.
Polarity
All magnets have two poles, one called the north pole and one called the south pole. Like poles repel and unlike polls attract (in
analogy to positive and negative charges in electrostatics). North and south poles always exist in pairs (there are no magnetic
monopoles in nature), so if one were to split a permanent magnet in half, two smaller magnets would be created, each with a north
pole and south pole.
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North and South Poles Always Come in Pairs: North and south poles always occur in pairs. Attempts to separate them result in
more pairs of poles. If we continue to split the magnet, we will eventually get down to an iron atom with a north pole and a south
pole—these, too, cannot be separated.
Making a Ferromagnet: An unmagnetized piece of iron is placed between two magnets, heated, and then cooled, or simply tapped
when cold. The iron becomes a permanent magnet with the poles aligned as shown: its south pole is adjacent to the north pole of
the original magnet, and its north pole is adjacent to the south pole of the original magnet. Note that there are attractive forces
between the magnets.
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Magnetic Field Lines
Magnetic field lines are useful for visually representing the strength and direction of the magnetic field.
learning objectives
Relate the strength of the magnetic field with the density of the magnetic field lines
Visualizing Magnetic Field Lines: Magnetic field lines are defined to have the direction that a small compass points when placed
at a location. (A) If small compasses are used to map the magnetic field around a bar magnet, they will point in the directions
shown: away from the north pole of the magnet, toward the south pole of the magnet (recall that Earth’s north magnetic pole is
really a south pole in terms of definitions of poles on a bar magnet. ) (B) Connecting the arrows gives continuous magnetic field
lines. The strength of the field is proportional to the closeness (or density) of the lines. (C) If the interior of the magnet could be
probed, the field lines would be found to form continuous closed loops.
Mapping the magnetic field of an object is simple in principle. First, measure the strength and direction of the magnetic field at a
large number of locations (or at every point in space). Then, mark each location with an arrow (called a vector ) pointing in the
direction of the local magnetic field with its magnitude proportional to the strength of the magnetic field (producing a vector field).
You can “connect” the arrows to form magnetic field lines. The direction of the magnetic field at any point is parallel to the
direction of nearby field lines, and the local density of field lines can be made proportional to its strength.
Magnetic field lines are like the contour lines (constant altitude) on a topographic map in that they represent something continuous,
and a different mapping scale would show more or fewer lines. An advantage of using magnetic field lines as a representation is
that many laws of magnetism (and electromagnetism) can be stated completely and concisely using simple concepts such as the
“number” of field lines through a surface. These concepts can be quickly translated to their mathematical form. For example, the
number of field lines through a given surface is the surface integral of the magnetic field.
Bar Magnet and Magnetic Field Lines: The direction of magnetic field lines represented by the alignment of iron filings
sprinkled on paper placed above a bar magnet.
Various phenomena have the effect of “displaying” magnetic field lines as though the field lines are physical phenomena. For
example, iron filings placed in a magnetic field line up to form lines that correspond to “field lines. ” Magnetic fields’ lines are also
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visually displayed in polar auroras, in which plasma particle dipole interactions create visible streaks of light that line up with the
local direction of Earth’s magnetic field.
Small compasses used to test a magnetic field will not disturb it. (This is analogous to the way we tested electric fields with a small
test charge. In both cases, the fields represent only the object creating them and not the probe testing them. ) Figure 15051 shows
how the magnetic field appears for a current loop and a long straight wire, as could be explored with small compasses. A small
compass placed in these fields will align itself parallel to the field line at its location, with its north pole pointing in the direction of
B. Note the symbols used for field into and out of the paper. We’ll explore the consequences of these various sources of magnetic
fields in further sections.
Mapping Magnetic Field Lines: Small compasses could be used to map the fields shown here. (A) The magnetic field of a
circular current loop is similar to that of a bar magnet. (B) A long and straight wire creates a field with magnetic field lines forming
circular loops. (C) When the wire is in the plane of the paper, the field is perpendicular to the paper. Note that the symbols used for
the field pointing inward (like the tail of an arrow) and the field pointing outward (like the tip of an arrow).
Extensive exploration of magnetic fields has revealed a number of hard-and-fast rules. We use magnetic field lines to represent the
field (the lines are a pictorial tool, not a physical entity in and of themselves). The properties of magnetic field lines can be
summarized by these rules:
1. The direction of the magnetic field is tangent to the field line at any point in space. A small compass will point in the direction
of the field line.
2. The strength of the field is proportional to the closeness of the lines. It is exactly proportional to the number of lines per unit
area perpendicular to the lines (called the areal density).
3. Magnetic field lines can never cross, meaning that the field is unique at any point in space.
4. Magnetic field lines are continuous, forming closed loops without beginning or end. They go from the north pole to the south
pole.
The last property is related to the fact that the north and south poles cannot be separated. It is a distinct difference from electric
field lines, which begin and end on the positive and negative charges. If magnetic monopoles existed, then magnetic field lines
would begin and end on them.
Geomagnetism
Earth’s magnetic field is caused by electric currents in the molten outer core and varies with time.
learning objectives
Explain the origin of the Earth’s magnetic field and its importance for the life on Earth
Geomagnetism
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surface of Earth, its magnetic field can be closely approximated by the field of a magnetic dipole positioned at the center of Earth
and tilted at an angle of about 10° with respect to the rotational axis of Earth.
Physical Origin
Earth’s magnetic field is mostly caused by electric currents in the liquid outer core, which is composed of highly conductive molten
iron. A magnetic field is generated by a feedback loop: Current loops generate magnetic fields (Ampère’s law); a changing
magnetic field generates an electric field (Faraday’s law); and the electric and magnetic fields exert a force on the charges that are
flowing in currents (the Lorentz force). These effects can be combined into a partial differential equation called the magnetic
induction equation:
∂B
2
= η ∇ B + ∇ × (u × B) (21.1.7)
∂t
In this equation u is the velocity of the fluid, B is the magnetic field, and eta is the magnetic diffusivity. The first term on the right
hand side of the induction equation is a diffusion term. If Earth’s dynamo shut off, the dipole part would disappear in a few tens of
thousands of years. The motion of the molten outer iron core is sustained by convection, or motion driven by buoyancy. The
temperature increases toward the center of Earth, and the higher temperature of the fluid lower down makes it buoyant. The
Coriolis effect, caused by the overall planetary rotation, tends to organize the flow into rolls aligned along the north-south polar
axis.
Origin of Earth’s Magnetic Field: A schematic illustrating the relationship between motion of conducting fluid, organized into
rolls by the Coriolis force, and the magnetic field the motion generates.
Electric currents induced in the ionosphere generate magnetic fields (ionospheric dynamo region). Such a field is always generated
near where the atmosphere is closest to the sun, causing daily alterations that can deflect surface magnetic fields by as much as one
degree. Typical daily variations of field strength are about 25 nanoteslas (nT), with variations over a few seconds of typically
around 1 nT.
Time Variations
The geomagnetic field changes on time scales from milliseconds to millions of years. Shorter time scales mostly arise from currents
in the ionosphere (ionospheric dynamo region) and magnetosphere, and some changes can be traced to geomagnetic storms or daily
variations in currents. Changes over time scales of a year or more mostly reflect changes in Earth’s interior, particularly the iron-
rich core. Frequently, Earth’s magnetosphere is hit by solar flares causing geomagnetic storms, provoking displays of aurorae. At
present, the overall geomagnetic field is becoming weaker; the present strong deterioration corresponds to a 10 to 15 percent
decline over the last 150 years and has accelerated in the past several years. Geomagnetic intensity has declined almost
continuously from a maximum 35 percent above the modern value achieved approximately 2,000 years ago. Earth’s magnetic
North Pole is drifting from northern Canada toward Siberia with a presently accelerating rate—10 km per year at the beginning of
the 20th century, up to 40 km per year in 2003, and since then has only accelerated.
Although Earth’s field is generally well approximated by a magnetic dipole with its axis near the rotational axis, there are
occasional dramatic events where the North and South geomagnetic poles trade places. These events are called geomagnetic
reversals. Evidence for these events can be found worldwide in basalts, sediment cores taken from the ocean floors, and seafloor
magnetic anomalies. Reversals occur at apparently random intervals ranging from less than 0.1 million years to as much as 50
million years. The most recent such event, called the Brunhes-Matuyama reversal, occurred about 780,000 years ago.
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Key Points
A wire carrying electric current will produce a magnetic field with closed field lines surrounding the wire.
Another version of the right hand rules can be used to determine the magnetic field direction from a current—point the thumb in
the direction of the current, and the fingers curl in the direction of the magnetic field loops created by it. See.
The Biot-Savart Law can be used to determine the magnetic field strength from a current segment. For the simple case of an
μ I
infinite straight current-carrying wire it is reduced to the form B = 0
2πr
.
A more fundamental law than the Biot-Savart law is Ampere ‘s Law, which relates magnetic field and current in a general way.
It is written in integral form as ∮ B ⋅ dl = μ I , where Ienc is the enclosed current and μ0 is a constant.
0 enc
A current-carrying wire feels a force in the presence of an external magnetic field. It is found to be F = Bilsinθ , where ℓ is the
length of the wire, i is the current, and θ is the angle between the current direction and the magnetic field.
Permanent magnets are objects made from magnetized material and produce continual magnetic fields. Everyday examples
include refrigerator magnets used to hold notes on a refrigerator door.
Materials that can be magnetized, which are also the ones that are strongly attracted to a magnet, are called ferromagnetic.
Examples of these materials include iron, nickel, and cobalt.
The counterexample to a permanent magnet is an electromagnet, which only becomes magnetized when an electric current
flows through it.
Magnets always have a north pole and a south pole, so if one were to split a permanent magnet in half, two smaller magnets
would be created, each with a north pole and south pole.
Permanent magnets are made from ferromagnetic materials that are exposed to a strong external magnetic field and heated to
align their internal microcrystalline structure, making them very hard to demagnetize.
The magnetic field direction is the same direction a compass needle points, which is tangent to the magnetic field line at any
given point.
The strength of the B-field is inversely proportional to the distance between field lines. It is exactly proportional to the number
of lines per unit area perpendicular to the lines.
A magnetic field line can never cross another field line. The magnetic field is unique at every point in space.
Magnetic field lines are continuous and unbroken, forming closed loops. Magnetic field lines are defined to begin on the north
pole of a magnet and terminate on the south pole.
Earth is largely protected from the solar wind, a stream of energetic charged particles emanating from the sun, by its magnetic
field. These particles would strip away the ozone layer, which protects Earth from harmful ultraviolet rays.
Earth’s magnetic field is generated by a feedback loop in the liquid outer core: Current loops generate magnetic fields; a
changing magnetic field generates an electric field; and the electric and magnetic fields exert a force on the charges that are
flowing in currents (the Lorentz force).
The geomagnetic field varies with time. Currents in the ionosphere and magnetosphere cause changes over short time scales,
while dramatic geomagnetic reversald (where the North and South poles switch locations) occur at apparently random intervals
ranging from 0.1 to 50 million years.
Key Terms
Biot-Savart Law: An equation that describes the magnetic field generated by an electric current. It relates the magnetic field to
the magnitude, direction, length, and proximity of the electric current. The law is valid in the magnetostatic approximation, and
is consistent with both Ampère’s circuital law and Gauss’s law for magnetism.
Ampere’s Law: An equation that relates magnetic fields to electric currents that produce them. Using Ampere’s law, one can
determine the magnetic field associated with a given current or current associated with a given magnetic field, providing there
is no time changing electric field present.
permanent magnet: A material, or piece of such material, which retains its magnetism even when not subjected to any external
magnetic fields.
ferromagnetic: Of a material, such as iron or nickel, that is easily magnetized.
electromagnet: A magnet which attracts metals only when electrically activated.
B-field: A synonym for the magnetic field.
magnetic field lines: A graphical representation of the magnitude and the direction of a magnetic field.
dynamo: A mechanism by which a celestial body such as Earth or a star generates a magnetic field over astronomical
timescales via a rotating, convecting, and electrically conducting fluid.
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21.2: Magnets
learning objectives
Identify two types of magnets
Ferromagnets
Only certain materials (e.g., iron, cobalt, nickel, and gadolinium) exhibit strong magnetic effects. These materials are called
ferromagnetic, after the Latin word ferrum (iron). A group of materials made from the alloys of the rare earth elements are also
used as strong and permanent magnets (neodymium is a common one). Other materials exhibit weak magnetic effects detectable
only with sensitive instruments. Not only do ferromagnetic materials respond strongly to magnets (the way iron is attracted to
magnets), they can also be magnetized themselves—that is, they can be induced to become magnetic or made into permanent
magnets.
When a magnet is brought near a previously unmagnetized ferromagnetic material, it causes local magnetization of the material
with unlike poles closest, as in. This results in the attraction of the previously unmagnetized material to the magnet as diagramed
in. When current produces a magnetic field on a microscopic scale, as illustrated in, the regions within the material called magnetic
domainsact like small bar magnets. Within domains, the poles of individual atoms are aligned, and each atom acts like a tiny bar
magnet. In an unmagnetized ferromagnetic object, domains are small and randomly oriented. In response to an external magnetic
field, the domains may grow to millimeter size, aligning themselves as shown in part (b) of the second figure. This induced
magnetization can become permanent if the material is heated and then cooled, or simply tapped in the presence of other magnets.
Current Produces a Magnetic Field: Current (I) through a wire produces a magnetic field (B). The field is oriented according to
the right-hand rule.
21.2.1 https://phys.libretexts.org/@go/page/15647
Unmagnetized to Magnetized Iron: (a) An unmagnetized piece of iron (or other ferromagnetic material) has randomly oriented
domains. (b) When magnetized by an external field, the domains show greater alignment, and some grow at the expense of others.
Individual atoms are aligned within domains; each atom acts like a tiny bar magnet.
Conversely, a permanent magnet can be demagnetized by hard blows or by heating it in the absence of another magnet. Increased
thermal motion at higher temperature can disrupt and randomize the orientation and the size of the domains. There is a well-defined
temperature for ferromagnetic materials called the Curie temperature, above which they cannot be magnetized. The Curie
temperature for iron is well above room temperature at 1043 K (770ºC). Several elements and alloys have Curie temperatures much
lower than room temperature, and are ferromagnetic only below those temperatures.
Electromagnets
In an electromagnet the magnetic field is produced by the flow of electric current. If the current disappears, the magnetic field is
turned off. Electromagnets are widely used as components of electrical devices, such as motors, generators, relays, loudspeakers,
hard disks, MRI machines, scientific instruments and magnetic separation equipment; they are also employed as industrial lifting
electromagnets for picking up and moving heavy iron objects like scrap iron.
An electric current flowing in a wire creates a magnetic field around the wire. To concentrate the magnetic field, the wire is wound
into a coil with many turns of wire lying side by side. The magnetic field from all the turns of wire passes through the center of the
coil creating a strong magnetic field there. The coil forming the shape of a straight tube (a helix) is called a solenoid, as shown in.
Much stronger magnetic fields can be produced if a “core” of ferromagnetic material (such as soft iron) is placed inside the coil.
Due to the high magnetic permeability μ of the ferromagnetic material, the ferromagnetic core increases the magnetic field to
thousands of times the strength of the field of the coil alone. This is called a ferromagnetic-core or iron-core electromagnet.
Electromagnet (Solenoid): A simple electromagnet consisting of a coil of insulated wire wrapped around an iron core. The
strength of magnetic field generated is proportional to the amount of current.
The direction of the magnetic field through a coil of wire can be likend to a form of the right-hand rule. If the fingers of the right
hand are curled around the coil in the direction of current flow (conventional current, flow of positive charge) through the
windings, the thumb points in the direction of the field inside the coil. The side of the magnet from which the field lines emerge is
defined as the north pole.The main advantage of an electromagnet over a permanent magnet is that the magnetic field can be
rapidly manipulated over a wide range by controlling the amount of electric current; a continuous supply of electrical energy is
required to maintain the field.
Key Points
Only certain materials, such as iron, cobalt, nickel, and gadolinium, exhibit strong magnetic effects. These materials are called
ferromagnetic. Ferromagnetic materials will respond strongly to magnets and can also be magnetized themselves.
Regions of uniform called magnetic domains are randomly oriented in unmagnetized ferromagnetic material, but may become
aligned under the influence of an external magnetic field. This process may become permanent if heated and cooled in the
presence of a magnetic field.
A ferromagnet will lose its magnetism if heated about its Curie temperature.
21.2.2 https://phys.libretexts.org/@go/page/15647
Electromagnets are a type of magnet in which the magnetic field is produced by the flow of current.
A strong electromagnet called a solenoid may be produced by wrapping wires into a coil and passing a current through them.
The magnetic field of all the turns of wire passes through the center of the coil, creating a strong magnetic field there.
Key Terms
magnetic domain: A region within a magnetic material which has uniform magnetization. This means that the individual
magnetic moments of the atoms are aligned with one another and they point in the same direction.
Curie temperature: The temperature above which a material will lose its magnetism.
solenoid: A coil of wire that acts as a magnet when an electric current flows through it.
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21.3: Magnetic Force on a Moving Electric Charge
Magnitude of the Magnetic Force
How does one magnet attracts another? The answer relies on the fact that all magnetism relies on current, the flow of charge.
Magnetic fields exert forces on moving charges, and so they exert forces on other magnets, all of which have moving charges.
The magnetic force on a moving charge is one of the most fundamental known. The magnetic force is as important as the
electrostatic or Coulomb force. Yet the magnetic force is more complex, in both the number of factors that affects it and in its
direction, than the relatively simple Coulomb force. The magnitude of the magnetic force F on a charge q moving at a speed v in a
magnetic field of strength B is given by:
F = qvB sin(θ) (21.3.1)
where θ is the angle between the directions of v and B . This formula is used to define the magnetic strength B in terms of the force
on a charged particle moving in a magnetic field. The SI unit for magnitude of the magnetic field strength is called the tesla (T) in
honor of the brilliant and eccentric inventor Nikola Tesla (1856–1943), who made great contributions to our understanding of
magnetic fields and their practical applications. To determine how the tesla relates to other SI units, we solve F = qvB sin(θ) for
B:
F
B= (21.3.2)
qvsin(θ)
Another smaller unit, called the gauss (G), where 1 G=10−4 T, is sometimes used. The strongest permanent magnets have fields
near 2 T; superconducting electromagnets may attain 10 T or more. The Earth’s magnetic field on its surface is only about 5×10−5
T, or 0.5 G.
The direction of the magnetic force F is perpendicular to the plane formed by v and B as determined by the right hand rule, which
is illustrated in Figure 1. It states that, to determine the direction of the magnetic force on a positive moving charge, you point the
thumb of the right hand in the direction of v , the fingers in the direction of B , and a perpendicular to the palm points in the
direction of F . One way to remember this is that there is one velocity, and so the thumb represents it. There are many field lines,
and so the fingers represent them. The force is in the direction you would push with your palm. The force on a negative charge is in
exactly the opposite direction to that on a positive charge.
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Right Hand Rule: Magnetic fields exert forces on moving charges. This force is one of the most basic known. The direction of the
magnetic force on a moving charge is perpendicular to the plane formed by v and B and follows right hand rule–1 (RHR-1) as
shown. The magnitude of the force is proportional to q, v, B, and the sine of the angle between v and B.
learning objectives
Apply the right hand rule to determine the direction of the magnetic force on a charge
Right Hand Rule: Magnetic fields exert forces on moving charges. This force is one of the most basic known. The direction of the
magnetic force on a moving charge is perpendicular to the plane formed by v and B and follows right hand rule–1 (RHR-1) as
shown. The magnitude of the force is proportional to q, v, B, and the sine of the angle between v and B.
One way to remember this is that there is one velocity, represented accordingly by the thumb. There are many field lines,
represented accordingly by the fingers. The force is in the direction you would push with your palm. The force on a negative charge
is in exactly the opposite direction to that on a positive charge. Because the force is always perpendicular to the velocity vector, a
pure magnetic field will not accelerate a charged particle in a single direction, however will produce circular or helical motion (a
concept explored in more detail in future sections). It is important to note that magnetic field will not exert a force on a static
electric charge. These two observations are in keeping with the rule that magnetic fields do no work.
Key Points
Magnetic fields exert forces on charged particles in motion.
The direction of the magnetic force F is perpendicular to the plane formed by v and B as determined by the right hand rule.
The SI unit for magnitude of the magnetic field strength is called the tesla (T), which is equivalent to one Newton per ampere-
meter. Sometimes the smaller unit gauss (10-4 T) is used instead.
When the expression for the magnetic force is combined with that for the electric force, the combined expression is known as
the Lorentz force.
21.3.2 https://phys.libretexts.org/@go/page/15648
When considering the motion of a charged particle in a magnetic field, the relevant vectors are the magnetic field B, the
velocity of the particle v, and the magnetic force exerted on the particle F. These vectors are all perpendicular to each other.
The right hand rule states that, to find the direction of the magnetic force on a positive moving charge, the thumb of the right
hand point in the direction of v, the fingers in the direction of B, and the force (F) is directed perpendicular to the right hand
palm.
The direction of the force F on a negative charge is in the opposite sense to that above (so pointed away from the back of your
hand).
Key Terms
Coulomb force: the electrostatic force between two charges, as described by Coulomb’s law
magnetic field: A condition in the space around a magnet or electric current in which there is a detectable magnetic force, and
where two magnetic poles are present.
tesla: In the International System of Units, the derived unit of magnetic flux density or magnetic inductivity. Symbol: T
right hand rule: Direction of angular velocity ω and angular momentum L in which the thumb of your right hand points when
you curl your fingers in the direction of rotation.
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21.4: Motion of a Charged Particle in a Magnetic Field
learning objectives
Compare the effects of the electric and the magnetic fields on the charged particle
Where F is the force vector, q is the charge, and E is the electric field vector. Note that the direction of F is identical to E in the case
of a positivist charge q, and in the opposite direction in the case of a negatively charged particle. This electric field may be
established by a larger charge, Q, acting on the smaller charge q over a distance r so that:
∣F∣ ∣ qQ ∣ |Q|
E =∣ ∣ = k∣ ∣ =k (21.4.2)
∣ q ∣ ∣ qr2 ∣ r
2
It should be emphasized that the electric force F acts parallel to the electric field E. The curl of the electric force is zero, i.e.:
∇×E = 0 (21.4.3)
A consequence of this is that the electric field may do work and a charge in a pure electric field will follow the tangent of an
electric field line.
In contrast, recall that the magnetic force on a charged particle is orthogonal to the magnetic field such that:
where B is the magnetic field vector, v is the velocity of the particle and θ is the angle between the magnetic field and the particle
velocity. The direction of F can be easily determined by the use of the right hand rule.
Right Hand Rule: Magnetic fields exert forces on moving charges. This force is one of the most basic known. The direction of the
magnetic force on a moving charge is perpendicular to the plane formed by v and B and follows right hand rule–1 (RHR-1) as
shown. The magnitude of the force is proportional to q, v, B, and the sine of the angle between v and B.
21.4.1 https://phys.libretexts.org/@go/page/15649
If the particle velocity happens to be aligned parallel to the magnetic field, or is zero, the magnetic force will be zero. This differs
from the case of an electric field, where the particle velocity has no bearing, on any given instant, on the magnitude or direction of
the electric force.
The angle dependence of the magnetic field also causes charged particles to move perpendicular to the magnetic field lines in a
circular or helical fashion, while a particle in an electric field will move in a straight line along an electric field line.
A further difference between magnetic and electric forces is that magnetic fields do not net work, since the particle motion is
circular and therefore ends up in the same place. We express this mathematically as:
W =∮ B ⋅ dr = 0 (21.4.5)
Lorentz Force
The Lorentz force is the combined force on a charged particle due both electric and magnetic fields, which are often considered
together for practical applications. If a particle of charge q moves with velocity v in the presence of an electric field E and a
magnetic field B, then it will experience a force:
F = q[E + vB sin θ] (21.4.6)
Electric Field Generated by Point Charges: The electric field surrounding three different point charges: (a) A positive charge; (b)
a negative charge of equal magnitude; (c) a larger negative charge.
If multiple charges are involved, field lines are generated on positive charges, and terminate on negative ones.
In the case of magnets, field lines are generated on the north pole (+) and terminate on the south pole (-) – see the below figure.
Magnetic ‘charges’, however, always come in pairs – there are no magnetic monopoles (isolated north or south poles). The curl of a
magnetic field generated by a conventional magnet is therefore always non zero. Charged particles will spiral around these field
lines, as long as the particles have some non-zero component of velocity directed perpendicular to the field lines.
21.4.2 https://phys.libretexts.org/@go/page/15649
Magnetic Pole Model: The magnetic pole model: two opposing poles, North (+) and South (−), separated by a distance d produce
an H-field (lines).
A magnetic field may also be generated by a current with the field lines envisioned as concentric circles around the current-
carrying wire.The magnetic force at any point in this case can be determined with the right hand rule, and will be perpendicular to
both the current and the magnetic field.
learning objectives
Identify conditions required for the particle to move in a straight line in the magnetic field
21.4.3 https://phys.libretexts.org/@go/page/15649
Zero Force When Velocity is Parallel to Magnetic Field: In the case above the magnetic force is zero because the velocity is
parallel to the magnetic field lines.
F = qvB sin θ (21.4.7)
If the magnetic field and the velocity are parallel (or antiparallel), then sinθ equals zero and there is no force. In this case a charged
particle can continue with straight-line motion even in a strong magnetic field. If is between 0 and 90 degrees, then the component
of v parallel to B remains unchanged.
Circular Motion
Since the magnetic force is always perpendicular to the velocity of a charged particle, the particle will undergo circular motion.
learning objectives
Describe conditions that lead to the circular motion of a charged particle in the magnetic field
Bubble Chamber: Trails of bubbles are produced by high-energy charged particles moving through the superheated liquid
hydrogen in this artist’s rendition of a bubble chamber. There is a strong magnetic field perpendicular to the page that causes the
curved paths of the particles. The radius of the path can be used to find the mass, charge, and energy of the particle.
So, does the magnetic force cause circular motion? Magnetic force is always perpendicular to velocity, so that it does no work on
the charged particle. The particle’s kinetic energy and speed thus remain constant. The direction of motion is affected, but not the
speed. This is typical of uniform circular motion. The simplest case occurs when a charged particle moves perpendicular to a
uniform B-field, such as shown in. (If this takes place in a vacuum, the magnetic field is the dominant factor determining the
motion. ) Here, the magnetic force (Lorentz force) supplies the centripetal force
21.4.4 https://phys.libretexts.org/@go/page/15649
Circular Motion of Charged Particle in Magnetic Field: A negatively charged particle moves in the plane of the page in a region
where the magnetic field is perpendicular into the page (represented by the small circles with x’s—like the tails of arrows). The
magnetic force is perpendicular to the velocity, and so velocity changes in direction but not magnitude. Uniform circular motion
results.
2
mv
Fc = (21.4.8)
r
Noting that
sin θ = 1 (21.4.9)
we see that
F = qvB (21.4.10)
The Lorentz magnetic force supplies the centripetal force, so these terms are equal:
2
mv
qvB = (21.4.11)
r
Here, r, called the gyroradius or cyclotron radius, is the radius of curvature of the path of a charged particle with mass m and charge
q, moving at a speed v perpendicular to a magnetic field of strength B. In other words, it is the radius of the circular motion of a
charged particle in the presence of a uniform magnetic field. If the velocity is not perpendicular to the magnetic field, then v is the
component of the velocity perpendicular to the field. The component of the velocity parallel to the field is unaffected, since the
magnetic force is zero for motion parallel to the field. We’ll explore the consequences of this case in a later section on spiral
motion.
A particle experiencing circular motion due to a uniform magnetic field is termed to be in a cyclotron resonance. The term comes
from the name of a cyclic particle accelerator called a cyclotron, showed in. The cyclotron frequency (or, equivalently,
gyrofrequency) is the number of cycles a particle completes around its circular circuit every second and can be found by solving for
v above and substituting in the circulation frequency so that
21.4.5 https://phys.libretexts.org/@go/page/15649
Cyclotron: A French cyclotron, produced in Zurich, Switzerland in 1937
v
f = (21.4.13)
2πr
becomes
qB
f = (21.4.14)
2πm
Helical Motion
Helical motion results when the velocity vector is not perpendicular to the magnetic field vector.
learning objectives
Describe conditions that lead to the helical motion of a charged particle in the magnetic field
Helical Motion
In the section on circular motion we described the motion of a charged particle with the magnetic field vector aligned perpendicular
to the velocity of the particle. In this case, the magnetic force is also perpendicular to the velocity (and the magnetic field vector, of
course) at any given moment resulting in circular motion. The speed and kinetic energy of the particle remain constant, but the
direction is altered at each instant by the perpendicular magnetic force. quickly reviews this situation in the case of a negatively
charged particle in a magnetic field directed into the page.
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Circular Motion of Charged Particle in Magnetic Field: A negatively charged particle moves in the plane of the page in a region
where the magnetic field is perpendicular into the page (represented by the small circles with x’s—like the tails of arrows). The
magnetic force is perpendicular to the velocity, and so velocity changes in direction but not magnitude. Uniform circular motion
results.
What if the velocity is not perpendicular to the magnetic field? Then we consider only the component of v that is perpendicular to
the field when making our calculations, so that the equations of motion become:
2
mv
⊥
Fc = (21.4.16)
r
The component of the velocity parallel to the field is unaffected, since the magnetic force is zero for motion parallel to the field.
This produces helical motion (i.e., spiral motion) rather than a circular motion.
shows how electrons not moving perpendicular to magnetic field lines follow the field lines. The component of velocity parallel to
the lines is unaffected, and so the charges spiral along the field lines. If field strength increases in the direction of motion, the field
will exert a force to slow the charges (and even reverse their direction), forming a kind of magnetic mirror.
Helical Motion and Magnetic Mirrors: When a charged particle moves along a magnetic field line into a region where the field
becomes stronger, the particle experiences a force that reduces the component of velocity parallel to the field. This force slows the
motion along the field line and here reverses it, forming a “magnetic mirror. “
The motion of charged particles in magnetic fields are related to such different things as the Aurora Borealis or Aurora Australis
(northern and southern lights) and particle accelerators. Charged particles approaching magnetic field lines may get trapped in
spiral orbits about the lines rather than crossing them, as seen above. Some cosmic rays, for example, follow the Earth’s magnetic
field lines, entering the atmosphere near the magnetic poles and causing the southern or northern lights through their ionization of
molecules in the atmosphere. Those particles that approach middle latitudes must cross magnetic field lines, and many are
prevented from penetrating the atmosphere. Cosmic rays are a component of background radiation; consequently, they give a
higher radiation dose at the poles than at the equator.
21.4.7 https://phys.libretexts.org/@go/page/15649
Charged Particles Spiral Along Earth’s Magnetic Field Lines: Energetic electrons and protons, components of cosmic rays,
from the Sun and deep outer space often follow the Earth’s magnetic field lines rather than cross them. (Recall that the Earth’s
north magnetic pole is really a south pole in terms of a bar magnet. )
learning objectives
Discuss application of mass spectrometers, movement of charged particles in a cyclotron, and how microwaves are
generated in the cavity magnetron
21.4.8 https://phys.libretexts.org/@go/page/15649
B
U~
~
Cyclotron Sketch: Sketch of a particle being accelerated in a cyclotron, and being ejected through a beamline.
Cyclotrons accelerate charged particle beams using a high frequency alternating voltage which is applied between two “D”-shaped
electrodes (also called “dees”). An additional static magnetic field is applied in perpendicular direction to the electrode plane,
enabling particles to re-encounter the accelerating voltage many times at the same phase. To achieve this, the voltage frequency
must match the particle’s cyclotron resonance frequency,
qB
f = (21.4.18)
2πm
with the relativistic mass m and its charge q. This frequency is given by equality of centripetal force and magnetic Lorentz force.
The particles, injected near the center of the magnetic field, increase their kinetic energy only when recirculating through the gap
between the electrodes; thus they travel outwards along a spiral path. Their radius will increase until the particles hit a target at the
perimeter of the vacuum chamber, or leave the cyclotron using a beam tube, enabling their use. The particles accelerated by the
cyclotron can be used in particle therapy to treat some types of cancer. Additionally, cyclotrons are a good source of high-energy
beams for nuclear physics experiments.
A synchrotron is an improvement upon the cyclotron in which the guiding magnetic field (bending the particles into a closed path)
is time-dependent, being synchronized to a particle beam of increasing kinetic energy. The synchrotron is one of the first
accelerator concepts that enable the construction of large-scale facilities, since bending, beam focusing and acceleration can be
separated into different components.
Cavity Magnetron
The cavity magnetron is a high-powered vacuum tube that generates microwaves using the interaction of a stream of electrons with
a magnetic field. All cavity magnetrons consist of a hot cathode with a high (continuous or pulsed) negative potential created by a
high-voltage, direct-current power supply. The cathode is built into the center of an evacuated, lobed, circular chamber. A magnetic
field parallel to the filament is imposed by a permanent magnet. The magnetic field causes the electrons, attracted to the (relatively)
positive outer part of the chamber, to spiral outward in a circular path, a consequence of the Lorentz force. Spaced around the rim
of the chamber are cylindrical cavities. The cavities are open along their length and connect the common cavity space. As electrons
sweep past these openings, they induce a resonant, high-frequency radio field in the cavity, which in turn causes the electrons to
bunch into groups.
21.4.9 https://phys.libretexts.org/@go/page/15649
Cavity Magnetron Diagram: A cross-sectional diagram of a resonant cavity magnetron. Magnetic lines of force are parallel to the
geometric axis of this structure.
The sizes of the cavities determine the resonant frequency, and thereby the frequency of emitted microwaves. The magnetron is a
self-oscillating device requiring no external elements other than a power supply. The magnetron has practical applications in radar,
heating (as the primary component of a microwave oven), and lighting.
Mass Spectrometry
Mass spectrometry is an analytical technique that measures the mass-to-charge ratio of charged particles. It is used for determining
masses of particles and determining the elemental composition of a sample or molecule.
Mass analyzers separate the ions according to their mass-to-charge ratio. The following two laws govern the dynamics of charged
particles in electric and magnetic fields in a vacuum:
F = Q(E + v × B) (Lorentz force) (21.4.19)
F = ma (21.4.20)
Equating the above expressions for the force applied to the ion yields:
(m/Q)a = E + v × B (21.4.21)
This differential equation along with initial conditions completely determines the motion of a charged particle in terms of m/Q.
There are many types of mass analyzers, using either static or dynamic fields, and magnetic or electric fields, but all operate
according to the above differential equation.
The following figure illustrates one type of mass spectrometer. The deflections of the particles are dependent on the mass-to-charge
ratio. In the case of isotopic carbon dioxide, each molecule has the same charge, but different masses. The mass spectrometer will
segregate the particles spatially allowing a detector to measure the mass-to-charge ratio of each particle. Since the charge is known,
the absolute mass can be determined trivially. The relative abundances can be inferred from counting the number of particles of
each given mass.
21.4.10 https://phys.libretexts.org/@go/page/15649
Mass Spectrometry: Schematics of a simple mass spectrometer with sector type mass analyzer. This one is for the measurement of
carbon dioxide isotope ratios (IRMS) as in the carbon-13 urea breath test.
Key Points
The force on a charged particle due to an electric field is directed parallel to the electric field vector in the case of a positive
charge, and anti-parallel in the case of a negative charge. It does not depend on the velocity of the particle.
In contrast, the magnetic force on a charge particle is orthogonal to the magnetic field vector, and depends on the velocity of the
particle. The right hand rule can be used to determine the direction of the force.
An electric field may do work on a charged particle, while a magnetic field does no work.
The Lorentz force is the combination of the electric and magnetic force, which are often considered together for practical
applications.
Electric field lines are generated on positive charges and terminate on negative ones. The field lines of an isolated charge are
directly radially outward. The electric field is tangent to these lines.
Magnetic field lines, in the case of a magnet, are generated at the north pole and terminate on a south pole. Magnetic poles do
not exist in isolation. Like in the case of electric field lines, the magnetic field is tangent to the field lines. Charged particles will
spiral around these field lines.
Newton’s first law of motion states that if an object experiences no net force, then its velocity is constant.
A particle with constant velocity will move along a straight line through space.
If a charged particle’s velocity is completely parallel to the magnetic field, the magnetic field will exert no force on the particle
and thus the velocity will remain constant.
In the case that the velocity vector is neither parallel nor perpendicular to the magnetic field, the component of the velocity
parallel to the field will remain constant.
The magnetic field does no work, so the kinetic energy and speed of a charged particle in a magnetic field remain constant.
The magnetic force, acting perpendicular to the velocity of the particle, will cause circular motion.
2
The centripetal force of the particle is provided by magnetic Lorentzian force so that qvB = mv
r
.
Solving for r above yields the gryoradius, or the radius of curvature of the path of a particle with charge q and mass m moving
in a magnetic field of strength B. The gryoradius is then given by Γ = mv
qB
.
The cyclotron frequency (or, equivalently, gyrofrequency) is the number of cycles a particle completes around its circular circuit
qB
every second and is given by f = 2πm
.
Previously, we have seen that circular motion results when the velocity of a charged particle is perpendicular to the magnetic
field. The speed and kinetic energy of the particle remain constant, but the direction is altered at each instant by the
perpendicular magnetic force.
If the velocity is not perpendicular to the magnetic field, we consider only the component of v that is perpendicular to the field
when making our calculations.
The component of the velocity parallel to the field is unaffected, since the magnetic force is zero for motion parallel to the field.
This produces helical motion.
21.4.11 https://phys.libretexts.org/@go/page/15649
Charges may spiral along field lines. If the strength of the magnetic field increases in the direction of motion, the field will exert
a force to slow the charges and even reverse their direction. This is known as a magnetic mirror.
A cyclotron is a type of particle accelerator in which charged particles accelerate outwards from the center along a spiral path.
The particles are held to a spiral trajectory by a static magnetic field and accelerated by a rapidly varying electric field.
The cavity magnetron is a high-powered vacuum tube that generates microwaves using the interaction of a stream of electrons
with a magnetic field. The magnetron has applications in radar, heating, and lighting.
Mass spectrometers measure the mass-to-charge ratio of charged particles through the use of electromagnetic fields to segregate
particles with different masses and/or charges. It can be used to determine the elemental composition of a molecule or sample.
Key Terms
orthogonal: Of two objects, at right angles; perpendicular to each other.
straight-line motion: motion that proceeds in a single direction.
gyroradius: The radius of the circular motion of a charged particle in the presence of a uniform magnetic field.
cyclotron frequency: The frequency of a charged particle moving perpendicular to the direction of a uniform magnetic field B
(constant magnitude and direction). Given by the equality of the centripetal force and magnetic Lorentz force.
helical motion: The motion that is produced when one component of the velocity is constant in magnitude and direction (i.e.,
straight-line motion) while the other component is constant in speed but uniformly varies in direction (i.e., circular motion). It is
the superposition of straight-line and circular motion.
magnetic mirror: A magnetic field configuration where the field strength changes when moving along a field line. The mirror
effect results in a tendency for charged particles to bounce back from the high field region.
cyclotron: An early particle accelerator in which charged particles were generated at a central source and accelerated spirally
outward through a fixed magnetic and alternating electric fields.
mass spectrometer: A device used in mass spectrometry to discover the mass composition of a given substance.
magnetron: A device in which electrons are made to resonate in a specially shaped chamber and thus produce microwave
radiation; used in radar, and in microwave ovens.
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21.5: Magnetic Fields, Magnetic Forces, and Conductors
learning objectives
Express Hall voltage for a a metal containing only one type of charge carriers
The Hall effect is the phenomenon in which a voltage difference (called the Hall voltage) is produced across an electrical
conductor, transverse to the conductor’s electric current when a magnetic field perpendicular to the conductor’s current is applied.
When a magnetic field is present that is not parallel to the motion of moving charges within a conductor, the charges experience the
Lorentz force. In the absence of such a field, the charges follow a roughly straight path, occasionally colliding with impurities.
In the presence of a magnetic field with a perpendicular component, the paths charges take becomes curved such that they
accumulate on one face of the material. On the other face, there is an excess of opposite charge remaining. Thus, an electric
potential is created so long as the charge flows. This opposes the magnetic force, eventually to the point of cancelation, resulting in
electron flow in a straight path.
Hall Effect for Electrons: Initially, the electrons are attracted by the magnetic force and follow the curved arrow. Eventually, when
electrons accumulate in excess on the left side and are in deficit on the right, an electric field ξy is created. This force becomes
strong enough to cancel out the magnetic force, so future electrons follow a straight (rather than curved) path.
For a metal containing only one type of charge carrier (electrons), the Hall voltage (VH) can be calculated as a factor of current (I),
magnetic field (B), thickness of the conductor plate (t), and charge carrier density (n) of the carrier electrons:
IB
VH = − (21.5.1)
net
The Hall effect is a rather ubiquitous phenomenon in physics, and appears not only in conductors, but semiconductors, ionized
gases, and in quantum spin among other applications.
learning objectives
Express equation used to calculate the magnetic force for an electrical wire exposed to a magnetic field
When an electrical wire is exposed to a magnet, the current in that wire will be affected by a magnetic field. The effect comes in the
form of a force. The expression for magnetic force on current can be found by summing the magnetic force on each of the many
individual charges that comprise the current. Since they all run in the same direction, the forces can be added.
21.5.1 https://phys.libretexts.org/@go/page/15650
Right Hand Rule: Used to determine direction of magnetic force.
The force (F) a magnetic field (B) exerts on an individual charge (q) traveling at drift velocity vd is:
F = qvdB sin θ (21.5.3)
In this instance, θ represents the angle between the magnetic field and the wire (magnetic force is typically calculated as a cross
product). If B is constant throughout a wire, and is 0 elsewhere, then for a wire with N charge carriers in its total length l, the total
magnetic force on the wire is:
F = NqvdB sin θ (21.5.4)
Given that N=nV, where n is the number of charge carriers per unit volume and V is volume of the wire, and that this volume is
calculated as the product of the circular cross-sectional area A and length (V=Al), yields the equation:
F = (nqAvd ) lB sin θ (21.5.5)
The terms in parentheses are equal to current (I), and thus the equation can be rewritten as:
F = IlB sin θ (21.5.6)
The direction of the magnetic force can be determined using the right hand rule, demonstrated in. The thumb is pointing in the
direction of the current, with the four other fingers parallel to the magnetic field. Curling the fingers reveals the direction of
magnetic force.
learning objectives
Identify the general quation for the torque on a loop of any shape
When a current travels in a loop that is exposed to a magnetic field, that field exerts torque on the loop. This principle is commonly
used in motors, in which the loop is connected to a shaft that rotates as a result of the torque. Thus, the electrical energy from the
current is converted to mechanical energy as the loop and shaft rotate, and this mechanical energy is then used to power another
device.
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Torque on a Current Loop: Electrical energy from the current is converted to mechanical energy as the loop and shaft rotate, and
this mechanical energy is then used to power another device.
In this model, the north and south poles of magnets are denoted by N and S, respectively. In the center is a rectangular wire loop of
length l and width w, carrying current I. The effect of magnetic field B on the current-carrying wire exerts torque τ.
To understand the torque, we must analyze the forces acting on each segment of the loop. Assuming a constant magnetic field, we
can conclude that the forces on the top and bottom parts of the loop are equal in magnitude and opposite in direction, and thus
produce no net force. Incidentally, those forces are vertical and thus parallel to the shaft.
However, as illustrated by (a) in the figure below, the equal but opposite forces produce a torque that acts clockwise.
Varying torque on a charged loop in a magnetic field: Maximum torque occurs in (b), when is 90 degrees. Minimum torque is 0,
and occurs in (c) when θ is 0 degrees. When loop rotates past =0, the torque reverses (d).
Given that torque is calculated from the equation:
τ = rF sin θ (21.5.7)
where F is force on the rotating object, r is the distance from the pivot point that the force is applied, and θ is the angle between r
and F, we can use the sum of two torques (the forces act on either side of the loop) to find the total torque:
w w
τ = F sin θ + F sin θ = wF sin θ (21.5.8)
2 2
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Note that the product of w and l is included in this equation; those terms can be replaced with area (A) of the rectangle. If another
shape of wire is used, its area can be inserted in the equation regardless of shape (whether circular, square, or otherwise).
Also note that this equation of torque is for a single turn. Torque increases proportionally according to number of turns (N). Thus,
the general equation for torque on a loop of any shape, of N turns, each of A area, carrying I current and exposed to a magnetic
field B is a value that fluctuates as the loop rotates, and can be calculated by:
learning objectives
Express the relationship between the strength of a magnetic field and a current running through a wire in a form of equation
Current running through a wire will produce both an electric field and a magnetic field. For a closed curve of length C, magnetic
field (B) is related to current (IC) as in Ampere’s Law, stated mathematically as:
∮ Bdℓ = μ0 IC (21.5.11)
C
Direction of magnetic field: The direction of the magnetic field can be determined by the right hand rule.
In this equation, dl represents the differential of length of wire in the curved wire, and μ0 is the permeability of free space. This can
be related to the Biot-Savart law. For a short, straight length of conductor (typically a wire) this law generally calculates partial
magnetic field (dB) as a function of current for an infinitesimally small segment of wire (dl) at a point r distance away from the
conductor:
μ0 Idl × r
dB = (21.5.12)
3
4π r
In this equation, the r vector can be written as r ̂ (the unit vector in direction of r), if the r3 term in the denominator is reduced to r2
(this is simply reducing like terms in a fraction). Integrating the previous differential equation, we find:
μ0 Idl × ^
r
B = ∮ (21.5.13)
2
4π C r
This relationship holds for constant current in a straight wire, in which magnetic field at a point due to all current elements
comprising the straight wire is the same. As illustrated in the direction of the magnetic field can be determined using the right hand
rule—pointing one’s thumb in the direction of current, the curl of one’s fingers indicates the direction of the magnetic field around
the straight wire.
21.5.4 https://phys.libretexts.org/@go/page/15650
learning objectives
Express the magnetic force felt by a pair of wires in a form of an equation
Parallel wires carrying current produce significant magnetic fields, which in turn produce significant forces on currents. The force
felt between the wires is used to define the the standard unit of current, know as an amphere.
In, the field (B1) that I1 creates can be calculated as a function of current and wire separation (r):
Magnetic fields and force exerted by parallel current-carrying wires.: Currents I1 and I2 flow in the same direction, separated
by a distance of r.
μ0 I1
B1 = (21.5.14)
2πr
The field B1 exerts a force on the wire containing I2. In the figure, this force is denoted as F2.
The force F2 exerts on wire 2 can be calculated as:
F2 = I2 lB1 sin θ (21.5.15)
Given that the field is uniform along and perpendicular to wire 2, sin θ = sin 90 derees = 1. Thus the force simplifies to: F2=I2lB1
According to Newton’s Third Law (F1=-F2), the forces on the two wires will be equal in magnitude and opposite in direction, so to
simply we can use F instead of F2. Given that wires are often very long, it’s often convenient to solve for force per unit length.
Rearranging the previous equation and using the definition of B1gives:
F μ0 I1 I2
= (21.5.16)
1 2πr
If the currents are in the same direction, the force attracts the wires. If the currents are in opposite directions, the force repels the
wires.
The force between current-carrying wires is used as part of the operational definition of the ampere. For parallel wires placed one
meter away from one another, each carrying one ampere, the force per meter is:
−7 2
F (4π ⋅ 10 T ⋅ m/A) (1A)
−7
= = 2 ⋅ 10 N/m (21.5.17)
l (2π)(1m)
Key Points
The Hall effect is the phenomenon in which a voltage difference (called the Hall voltage) is produced across an electrical
conductor that is transverse to the conductor’s electric current when a magnetic field perpendicular to the conductor’s current is
applied.
Moving charges in a wire will change trajectory in the presence of a magnetic field, “bending” toward it. Thus, those charges
accumulate on one face of the material. On the other face, there is left an excess of opposite charge. Thus, an electric potential
21.5.5 https://phys.libretexts.org/@go/page/15650
is created.
VH = −
IB
net
is the formula for Hall voltage (VH). It is a factor of current (I), magnetic field (B), thickness of the conductor plate
(t), and charge carrier density (n) of the carrier electrons.
Magnetic force on current can be found by summing the magnetic force on each of the individual charges that make this
current.
For a wire exposed to a magnetic field, τ = NIAB sin θ describes the relationship between magnetic force (F), current (I),
length of wire (l), magnetic field (B), and angle between field and wire (θ).
The direction of the magnetic force can be determined using the right hand rule, as in fig .
τ = NIAB sin θ can be used to calculate torque (ττ) a loop of N turns and A area, carrying I current feels in the presence of a
magnetic field B.
Although the forces acting upon the loop are equal and opposite, they both act to rotate the loop in the same direction.
Torque experienced is independent of the loop’s shape. What matters is the area of the loop.
Ampere ‘s Law states that for a closed curve of length C, magnetic field (B) is related to current (IC): ∮ Bdℓ = μ I . In this
C
0 C
equation, dl represents the differential of length of wire in the curved wire, and μ0 is the permeability of free space.
μ
Ampere’s Law can be related to the Biot-Savart law, which holds for a short, straight length of conductor: dB = 0
4π
. In
Id×lr
r3
this equation, partial magnetic field (dB) is expressed as a function of current for an infinitesimally small segment of wire (dl)
at a point r distance away from the conductor.
After integrating, the direction of the magnetic field according to the Biot-Savart Law can be determined using the right hand
rule.
The field (B1) that that current (I1) from a wire creates can be calculated as a function of current and wire separation (r):
μ I1
B1 =
0
2πr
μ0 is a constant.
F = IlB sin θ describes the magnetic force felt by a pair of wires. If they are parallel the equation is simplified as the sine
function is 1.
The force felt between two parallel conductive wires is used to define the ampere —the standard unit of current.
Key Terms
elementary charge: The electric charge on a single proton.
transverse: Not tangent, so that a nondegenerate angle is formed between the two things intersecting.
drift velocity: The average velocity of the free charges in a conductor.
magnetic field: A condition in the space around a magnet or electric current in which there is a detectable magnetic force, and
where two magnetic poles are present.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
electric field: A region of space around a charged particle, or between two voltages; it exerts a force on charged objects in its
vicinity.
ampere: A unit of electrical current; the standard base unit in the International System of Units. Abbreviation: amp. Symbol: A.
current: The time rate of flow of electric charge.
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21.6: Applications of Magnetism
learning objectives
Explain why magnetic field is utilized in mass spectrometers
Mass Spectrometry
Mass spectrometry (MS) is the art of displaying the spectra (singular spectrum) of the masses of a sample of material. It is used for
determining the elemental composition of a sample, and properties of the particles and molecules (the chemical structures of
molecules, such as peptides and other chemical compounds).
Schematic of Mass Spectrometer: Schematics of a simple mass spectrometer with sector type mass analyzer. This one is for the
measurement of carbon dioxide isotope ratios as in the carbon-13urea breath test.
Mass spectrometers, as diagramed in, separate compounds based on a property known as the mass-to-charge ratio. The sample to
be identified is first ionized, and then passed through some form of magnetic field. Based on parameters, such as how long it takes
the molecule to travel a certain distance or the amount of deflection caused by the field, a mass can be calculated for the ion.
How it Works
First, the sample undergoes vaporization by intense heat. The gaseous components of the sample are each ionized (turned into ions)
with the same quantity of charge. The ions are then grouped by applying a similar magnetic force to each. Since the acceleration of
a charge is dependent on the mass and strength of the charge, a lighter mass-to-charge ratio will not travel as far as a high mass-to-
charge ratio, allowing for comparison of the physical properties of different particles. The groupings are detected by some
quantitative signal from of ongoing probing. The detector differentiates the ions based on how far they curve in the magnetic field.
The signal is processed into the spectra (singular of spectrum) of the masses of the particles of that sample. The elements or
molecules are uniquely identified by correlating known masses by the identified masses.
Ferromagnetism
Ferromagnetism is the property of certain materials that enables them to form magnets and be attracted to magnets.
learning objectives
Relate ferromagnetism with the electron configuration
Ferromagnetism is the basic mechanism by which certain materials (such as iron) form permanent magnets, or are attracted to
magnets. In physics, several different types of magnetism are distinguished. Ferromagnetism is the strongest type—it is the only
type that creates forces strong enough to be felt, and is responsible for the common phenomena of magnetism encountered in
everyday life.
When a magnet comes into close proximity with a previously unmagnetized ferromagnetic material, it causes local magnetization
of the material with closest unlike poles. Regions within the material (called domains) act like small bar magnets. Within domains,
the poles of individual atoms are aligned. Each atom acts like a tiny bar magnet. Domains are small and randomly oriented in an
unmagnetized ferromagnetic object. In response to an external magnetic field, the domains may grow to millimeter size, aligning
themselves. This induced magnetization can be made permanent if the material is heated and then cooled, or simply tapped in the
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presence of other magnets, as shown in. Permanent magnets (materials that can be magnetized by an external magnetic field and
remain magnetized after the external field is removed) are ferromagnetic, as are other materials that are noticeably attracted to
them.
Unmagnetized to Magnetized Iron: (a) An unmagnetized piece of iron (or other ferromagnetic material) has randomly oriented
domains. (b) When magnetized by an external field, the domains show greater alignment, and some grow at the expense of others.
Individual atoms are aligned within domains; each atom acts like a tiny bar magnet.
Ferromagnetism arises from the fundamental property of an electron; it also carries charge to have a dipole moment. Thus, an
electron itself behaves as a tiny magnet. This dipole moment comes from the more fundamental property of the electron—its
quantum mechanical spin. The quantum mechanical nature of this spin limits the electron to only two states: with the magnetic
field pointing either “up” or “down” (for any choice of up and down). When these tiny magnetic dipoles are aligned in the same
direction, their individual magnetic fields combine to create a measurable macroscopic field.
However, in materials with a filled electron shell, the total dipole moment of the electrons is zero, as the spins are in up/down pairs.
Only atoms with partially filled shells (i.e., unpaired spins) can have a net magnetic moment. Thus ferromagnetism only occurs in
materials with partially filled shells. (According to Hund’s rules, the first few electrons in a shell tend to have the same spin,
thereby increasing the total dipole moment. )
Accordingly, only certain materials (such as iron, cobalt, nickel, and gadolinium) exhibit strong magnetic effects. Such materials
are called ferromagnetic, after the Latin word for iron, ferrum. A group of materials made from the alloys of the rare earth elements
are also used as strong and permanent magnets (a popular one is neodymium). Other materials exhibit weak magnetic effects,
detectable only with sensitive instruments. Not only do ferromagnetic materials respond strongly to magnets (the way iron is
attracted to magnets), they can also be magnetized themselves—that is, they can be induced to be magnetic or made into permanent
magnets.
Ferromagnetism is very important in industry and modern technology, and is the basis for many electrical and electromechanical
devices such as: electromagnets, electric motors, generators, transformers, and magnetic storage (e.g., tape recorders and hard
disks).
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Refrigerator Magnets: Different magnets attached to the doors of a refrigerator.
learning objectives
Describe properties of diamagnetic and paramagnetic materials
Paramagnetism
Paramagnetism is a form of magnetism whereby the paramagnetic material is only attracted when in the presence of an externally
applied magnetic field. Paramagnetic materials have a relative magnetic permeability greater or equal to unity (i.e., a positive
magnetic susceptibility) and hence are attracted to magnetic fields. The magnetic moment induced by the applied field is linear in
the field strength; it is also rather weak.
Constituent atoms or molecules of paramagnetic materials have permanent magnetic moments ( dipoles ), even in the absence of an
applied field. Generally, the permanent moment is caused by the spin of unpaired electrons in atomic or molecular electron orbitals.
In pure paramagnetism, the dipoles do not interact with each other and are randomly oriented in the absence of an external field due
to thermal agitation; this results in a zero net magnetic moment. When a magnetic field is applied, the dipoles will tend to align
with the applied field, resulting in a net magnetic moment in the direction of the applied field.
Paramagnetic materials have a small, positive susceptibility to magnetic fields. These materials are slightly attracted by a magnetic
field and the material does not retain the magnetic properties when the external field is removed, as illustrated in. Paramagnetic
properties are due to the presence of some unpaired electrons, and from the realignment of the electron paths caused by the external
magnetic field.
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Paramagnetic Materials and Electric Fields: Orientation in paramagnetic material when electric field is applied (right image)
and removed (left image).
Paramagnetic materials include magnesium, molybdenum, lithium and tantalum. Unlike ferromagnets, paramagnets do not retain
any magnetization in the absence of an externally applied magnetic field, because thermal motion randomizes the spin orientations
responsible for magnetism. Some paramagnetic materials retain spin disorder at absolute zero (meaning they are paramagnetic in
the ground state). Thus the total magnetization drops to zero when the applied field is removed. Even in the presence of the field
there is only a small induced magnetization because only a small fraction of the spins will be oriented by the field.
Diamagnetism
Diamagnetism is the property of an object or material that causes it to create a magnetic field in opposition to an externally applied
magnetic field. Thus, unlike paramagnets, diamagnets are repelled by magnetic fields, which can lead to its unusual effects, such as
levitation of diamagnetic material when located above powerful magnet (as shown in ).
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Solenoids, Current Loops, and Electromagnets
Solenoids are loops of wire around a metallic core, and can be used to create controlled magnetic fields.
learning objectives
Explain why solenoids are important and their structure
Solenoids
A solenoid is a coil wound into a tightly packed helix. In physics, the term solenoid refers to a long, thin loop of wire, often
wrapped around a metallic core; it produces a magnetic field when an electric current is passed through it. Solenoids are important
because they can create controlled magnetic fields and can be used as electromagnets. The term solenoid refers specifically to a coil
designed to produce a uniform magnetic field in a volume of space (in which some experiment might be performed).
Uses of Electromagnets: An electromagnet induces regions of permanent magnetism on a floppy disk coated with a ferromagnetic
material. The information stored here is digital (a region is either magnetic or not); in other applications, it can be analog (with a
varying strength), such as on audiotapes.
Combining a ferromagnet with an electromagnet can produce particularly strong magnetic effects. Whenever strong magnetic
effects are needed (such as lifting scrap metal, or in particle accelerators) electromagnets are enhanced by ferromagnetic materials.
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An electromagnet creates magnetism with an electric current. In later sections we explore this more quantitatively, finding the
strength and direction of magnetic fields created by various currents. Currents, including those associated with other
submicroscopic particles like protons, allow us to explain ferromagnetism and all other magnetic effects. Ferromagnetism, for
example, results from an internal cooperative alignment of electron spins, possible in some materials but not in others.
Key Points
By vaporizing a material into ions, mass spectrometers can tell what elements make up that material.
When a magnetic field is applied to ions, they deflect or accelerate various amounts depending on their electric charge to mass
ratio.
By observing how much the ions deflect, the mass spectrometer can identify the charge to mass ratio of each ion and thus which
element it is.
Ferromagnetism is the basic mechanism by which certain materials (such as iron) form permanent magnets, or are attracted to
magnets.
When a ferromagnetic material is brought close to a magnet, the poles of the individual atoms of the material align along the
magnetic field lines. If made permanent, this alignment can create a permanent magnet.
Ferromagnetism only occurs in materials with partially filled electron shells.
Paramagnets act like magnets while in the presence of an externally applied magnetic field.
Diamagnets create a magnetic field in opposition to an externally applied magnetic field. Thus, they repulse magnets.
Diamagnetism is a property of all materials and always makes a weak contribution to the material’s response to a magnetic
field. However, for materials that show ferromagnetism or paramagnetism, the diamagnetic contribution becomes negligible.
Loops of wire in a magnetic field create current.
Electromagnetism is the use of electric current to make magnets. Electromagnets are temporary magnets which keep their
magnetic properties only when current is passing through them.
Solenoids and electromagnets have many uses in physics and engineering because they allow the control of magnetic fields.
Key Terms
electric charge: A quantum number that determines the electromagnetic interactions of some subatomic particles; by
convention, the electron has an electric charge of -1 and the proton +1, and quarks have fractional charge.
mass: The quantity of matter which a body contains, irrespective of its bulk or volume. It is one of four fundamental properties
of matter. It is measured in kilograms in the SI system of measurement.
magnetic field: A condition in the space around a magnet or electric current in which there is a detectable magnetic force, and
where two magnetic poles are present.
spin: A quantum angular momentum associated with subatomic particles; it also creates a magnetic moment.
dipole moment: The vector product of the charge on either pole of a dipole and the distance separating them.
electron shell: The collective states of all electrons in an atom having the same principal quantum number (visualized as an
orbit in which the electrons move).
ferromagnetism: The phenomenon whereby certain substances can become permanent magnets when subjected to a magnetic
field.
paramagnetism: The tendency of magnetic dipoles to align with an external magnetic field; materials that exhibit this tendency
become temporary magnets.
diamagnetism: A weak form of magnetism that is only observed in the presence of an external magnetic field; due to an
induced magnetic field in an opposite direction.
ferromagnetic: Of a material, such as iron or nickel, that is easily magnetized.
magnetic field: A condition in the space around a magnet or electric current in which there is a detectable magnetic force, and
where two magnetic poles are present.
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CHAPTER OVERVIEW
22: Induction, AC Circuits, and Electrical Technologies is shared under a not declared license and was authored, remixed, and/or curated by
LibreTexts.
1
22.1: Magnetic Flux, Induction, and Faraday’s Law
learning objectives
Explain the relationship between the magnetic field and the electromotive force
Induced EMF
The apparatus used by Faraday to demonstrate that magnetic fields can create currents is illustrated in the following figure. When the switch is closed, a
magnetic field is produced in the coil on the top part of the iron ring and transmitted (or guided) to the coil on the bottom part of the ring. The galvanometer is
used to detect any current induced in a separate coil on the bottom.
Faraday’s Apparatus: This is Faraday’s apparatus for demonstrating that a magnetic field can produce a current. A change in the field produced by the top
coil induces an EMF and, hence, a current in the bottom coil. When the switch is opened and closed, the galvanometer registers currents in opposite directions.
No current flows through the galvanometer when the switch remains closed or open.
It was found that each time the switch is closed, the galvanometer detects a current in one direction in the coil on the bottom. Each time the switch is opened,
the galvanometer detects a current in the opposite direction. Interestingly, if the switch remains closed or open for any length of time, there is no current
through the galvanometer. Closing and opening the switch induces the current. It is the change in magnetic field that creates the current. More basic than the
current that flows is the electromotive force (EMF) that causes it. The current is a result of an EMF induced by a changing magnetic field, whether or not there
is a path for current to flow.
Magnetic Flux
The magnetic flux (often denoted Φ or ΦB) through a surface is the component of the magnetic field passing through that surface. The magnetic flux through
some surface is proportional to the number of field lines passing through that surface. The magnetic flux passing through a surface of vector area A is
ΦB = B ⋅ A = BA cos θ (22.1.1)
where B is the magnitude of the magnetic field (having the unit of Tesla, T), A is the area of the surface, and θ is the angle between the magnetic field lines and
the normal (perpendicular) to A.
For a varying magnetic field, we first consider the magnetic flux dΦBdΦB through an infinitesimal area element dA, where we may consider the field to be
constant:
Varying Magnetic Field: Each point on a surface is associated with a direction, called the surface normal; the magnetic flux through a point is then the
component of the magnetic field along this normal direction.
\mathrm { d } \Phi _ { \mathrm { B } } = \mathbf { B } \cdot \mathrm { d } \mathbf { A }
A generic surface, A, can then be broken into infinitesimal elements and the total magnetic flux through the surface is then the surface integral
ΦB = ∬ B ⋅ dA (22.1.2)
A
Δt
, when flux changes by Δ in a time Δt.
learning objectives
Express the Faraday’s law of induction in a form of equation
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Faraday’s Law of Induction
Faraday’s law of induction is a basic law of electromagnetism that predicts how a magnetic field will interact with an electric circuit to produce an
electromotive force (EMF). It is the fundamental operating principle of transformers, inductors, and many types of electrical motors, generators, and solenoids.
Faraday’s experiments showed that the EMF induced by a change in magnetic flux depends on only a few factors. First, EMF is directly proportional to the
change in flux Δ. Second, EMF is greatest when the change in time Δt is smallest—that is, EMF is inversely proportional to Δt. Finally, if a coil has N turns, an
EMF will be produced that is N times greater than for a single coil, so that EMF is directly proportional to N. The equation for the EMF induced by a change in
magnetic flux is
ΔΦ
EMF = −N (22.1.3)
Δt
This relationship is known as Faraday’s law of induction. The units for EMF are volts, as is usual.
Lenz’ Law
The minus sign in Faraday’s law of induction is very important. The minus means that the EMF creates a current I and magnetic field B that oppose the change
in flux Δthis is known as Lenz’ law. The direction (given by the minus sign) of the EMF is so important that it is called Lenz’ law after the Russian Heinrich
Lenz (1804–1865), who, like Faraday and Henry, independently investigated aspects of induction. Faraday was aware of the direction, but Lenz stated it, so he
is credited for its discovery.
Lenz’ Law: (a) When this bar magnet is thrust into the coil, the strength of the magnetic field increases in the coil. The current induced in the coil creates
another field, in the opposite direction of the bar magnet’s to oppose the increase. This is one aspect of Lenz’s law—induction opposes any change in flux. (b)
and (c) are two other situations. Verify for yourself that the direction of the induced Bcoil shown indeed opposes the change in flux and that the current
direction shown is consistent with the right hand rule.
Energy Conservation
Lenz’ law is a manifestation of the conservation of energy. The induced EMF produces a current that opposes the change in flux, because a change in flux
means a change in energy. Energy can enter or leave, but not instantaneously. Lenz’ law is a consequence. As the change begins, the law says induction
opposes and, thus, slows the change. In fact, if the induced EMF were in the same direction as the change in flux, there would be a positive feedback that
would give us free energy from no apparent source—conservation of energy would be violated.
Motional EMF
Motion in a magnetic field that is stationary relative to the Earth induces motional EMF (electromotive force).
learning objectives
Identify process that induces motional electromotive force
As seen in previous Atoms, any change in magnetic flux induces an electromotive force (EMF) opposing that change—a process known as induction. Motion
is one of the major causes of induction. For example, a magnet moved toward a coil induces an EMF, and a coil moved toward a magnet produces a similar
EMF. In this Atom, we concentrate on motion in a magnetic field that is stationary relative to the Earth, producing what is loosely called motional EMF.
Motional EMF
Consider the situation shown in. A rod is moved at a speed v along a pair of conducting rails separated by a distance ℓ in a uniform magnetic field B. The rails
are stationary relative to B, and are connected to a stationary resistor R (the resistor could be anything from a light bulb to a voltmeter). Consider the area
enclosed by the moving rod, rails and resistor. B is perpendicular to this area, and the area is increasing as the rod moves. Thus the magnetic flux enclosed by
the rails, rod and resistor is increasing. When flux changes, an EMF is induced according to Faraday’s law of induction.
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Motional EMF: (a) A motional emf=Bℓv is induced between the rails when this rod moves to the right in the uniform magnetic field. The magnetic field B is
into the page, perpendicular to the moving rod and rails and, hence, to the area enclosed by them. (b) Lenz’s law gives the directions of the induced field and
current, and the polarity of the induced emf. Since the flux is increasing, the induced field is in the opposite direction, or out of the page. Right hand rule gives
the current direction shown, and the polarity of the rod will drive such a current.
To find the magnitude of EMF induced along the moving rod, we use Faraday’s law of induction without the sign:
ΔΦ
EMF = N (22.1.4)
Δt
In this equation, N=1 and the flux Φ=BAcosθ. We have θ=0º and cosθ=1, since B is perpendicular to A. Now Δ = Δ(BA) = BΔA , since B is uniform. Note
that the area swept out by the rod is ΔA = ℓx . Entering these quantities into the expression for EMF yields:
BΔA lΔx
EMF = =B = Blv (22.1.5)
Δt Δt
To find the direction of the induced field, the direction of the current, and the polarity of the induced EMF we apply Lenz’ law, as explained in Faraday’s Law
of Induction: Lenz’ Law. As seen in Fig 1 (b), F lux is increasing, since the area enclosed is increasing. Thus the induced field must oppose the existing one
and be out of the page. (The right hand rule requires that I be counterclockwise, which in turn means the top of the rod is positive, as shown. )
learning objectives
Explain the relationship between the motional electromotive force, eddy currents, and magnetic damping
Back EMF
Motors and generators are very similar. (Read our Atoms on “Electric Generators” and “Electric Motors. “) Generators convert mechanical energy into
electrical energy, whereas motors convert electrical energy into mechanical energy. Furthermore, motors and generators have the same construction. When the
coil of a motor is turned, magnetic flux changes, and an electromotive force (EMF), consistent with Faraday’s law of induction, is induced. The motor thus acts
as a generator whenever its coil rotates. This will happen whether the shaft is turned by an external input, like a belt drive, or by the action of the motor itself.
That is, when a motor is doing work and its shaft is turning, an EMF is generated. Lenz’ law tells us the induced EMF opposes any change, so that the input
EMF that powers the motor will be opposed by the motor’s self-generated EMF, called the back EMF of the motor.
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Eddy Current
As discussed in “Motional EMF,” motional EMF is induced when a conductor moves in a magnetic field or when a magnetic field moves relative to a
conductor. If motional EMF can cause a current loop in the conductor, we refer to that current as an eddy current. Eddy currents can produce significant drag,
called magnetic damping, on the motion involved.
Consider the apparatus shown in, which swings a pendulum bob between the poles of a strong magnet. If the bob is metal, there is significant drag on the bob
as it enters and leaves the field, quickly damping the motion. If, however, the bob is a slotted metal plate, as shown in (b), there is a much smaller effect due to
the magnet. There is no discernible effect on a bob made of an insulator.
Device for Exploring Eddy Currents and Magnetic Damping: A common physics demonstration device for exploring eddy currents and magnetic damping.
(a) The motion of a metal pendulum bob swinging between the poles of a magnet is quickly damped by the action of eddy currents. (b) There is little effect on
the motion of a slotted metal bob, implying that eddy currents are made less effective. (c) There is also no magnetic damping on a nonconducting bob, since the
eddy currents are extremely small.
shows what happens to the metal plate as it enters and leaves the magnetic field. In both cases, it experiences a force opposing its motion. As it enters from the
left, flux increases, and so an eddy current is set up (Faraday’s law) in the counterclockwise direction (Lenz’ law), as shown. Only the right-hand side of the
current loop is in the field, so that there is an unopposed force on it to the left (right hand rule). When the metal plate is completely inside the field, there is no
eddy current if the field is uniform, since the flux remains constant in this region. But when the plate leaves the field on the right, flux decreases, causing an
eddy current in the clockwise direction that, again, experiences a force to the left, further slowing the motion. A similar analysis of what happens when the
plate swings from the right toward the left shows that its motion is also damped when entering and leaving the field.
Conducting Plate Passing Between the Poles of a Magnet: A more detailed look at the conducting plate passing between the poles of a magnet. As it enters
and leaves the field, the change in flux produces an eddy current. Magnetic force on the current loop opposes the motion. There is no current and no magnetic
drag when the plate is completely inside the uniform field.
When a slotted metal plate enters the field, as shown in, an EMF is induced by the change in flux, but it is less effective because the slots limit the size of the
current loops. Moreover, adjacent loops have currents in opposite directions, and their effects cancel. When an insulating material is used, the eddy current is
extremely small, and so magnetic damping on insulators is negligible. If eddy currents are to be avoided in conductors, then they can be slotted or constructed
of thin layers of conducting material separated by insulating sheets.
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Eddy Currents Induced in a Slotted Metal Plate: Eddy currents induced in a slotted metal plate entering a magnetic field form small loops, and the forces on
them tend to cancel, thereby making magnetic drag almost zero.
learning objectives
Describe the relationship between the changing magnetic field and an electric field
We have studied Faraday’s law of induction in previous atoms. We learned the relationship between induced electromotive force (EMF) and magnetic flux. In a
dΦB ∂ΦB
nutshell, the law states that changing magnetic field( dt
) produces an electric field (ε ), Faraday’s law of induction is expressed as ε =−
∂t
, where ε is
dΦB
induced EMF and dt
is magnetic flux. (“N” is dropped from our previous expression. The number of turns of coil is included can be incorporated in the
magnetic flux, so the factor is optional. ) Faraday’s law of induction is a basic law of electromagnetism that predicts how a magnetic field will interact with an
electric circuit to produce an electromotive force (EMF). In this Atom, we will learn about an alternative mathematical expression of the law.
Faraday’s Experiment: Faraday’s experiment showing induction between coils of wire: The liquid battery (right) provides a current which flows through the
small coil (A), creating a magnetic field. When the coils are stationary, no current is induced. But when the small coil is moved in or out of the large coil (B),
the magnetic flux through the large coil changes, inducing a current which is detected by the galvanometer (G).
current is a source of electromotive force. (EMF) The definition is mathematically ε = ∮ E⃗ ⋅ ds ⃗ , where the integral is evaluated over a closed loop C.
C
∂t
⃗ ⃗
(∫ B ⋅ dA) . Using the Stokes’ theorem in vector calculus, the left hand side is
⃗
∮
C
⃗ ⃗
E ⋅ ds ⃗ = ∫ (∇ × E) ⋅ dA
S
⃗
. Also, note that in the right hand side ∂
∂t
⃗ ⃗
(∫ B ⋅ dA) = ∫
∂B
∂t
⋅ dA
⃗
. Therefore, we get an alternative form of the Faraday’s law of
⃗
induction: ∇ × E⃗ = − .This is also called a differential form of the Faraday’s law. It is one of the four equations in Maxwell’s equations, governing all
∂B
∂t
electromagnetic phenomena.
Electric Generators
Electric generators convert mechanical energy to electrical energy; they induce an EMF by rotating a coil in a magnetic field.
22.1.5 https://phys.libretexts.org/@go/page/15654
learning objectives
Explain how an electromotive force is induced in electric generators
Electric generators are devices that convert mechanical energy to electrical energy. They induce an electromotive force (EMF) by rotating a coil in a magnetic
field. It is a device that converts mechanical energy to electrical energy. A generator forces electric charge (usually carried by electrons) to flow through an
external electrical circuit. Possible sources of mechanical energy include: a reciprocating or turbine steam engine, water falling through a turbine or
waterwheel, an internal combustion engine, a wind turbine, a hand crank, compressed air, or any other source of mechanical energy. Generators supply almost
all of the power for the electric power grids which provide most of the world’s electric power.
Basic Setup
Consider the setup shown in. Charges in the wires of the loop experience the magnetic force because they are moving in a magnetic field. Charges in the
vertical wires experience forces parallel to the wire, causing currents. However, those in the top and bottom segments feel a force perpendicular to the wire;
this force does not cause a current. We can thus find the induced EMF by considering only the side wires. Motional EMF is given to be EMF=Bℓv, where the
velocity v is perpendicular to the magnetic field B (see our Atom on “Motional EMF”). Here, the velocity is at an angle θ with B, so that its component
perpendicular to B is vsinθ.
Diagram of an Electric Generator: A generator with a single rectangular coil rotated at constant angular velocity in a uniform magnetic field produces an emf
that varies sinusoidally in time. Note the generator is similar to a motor, except the shaft is rotated to produce a current rather than the other way around.
Thus in this case the EMF induced on each side is EMF=Bℓvsinθ, and they are in the same direction. The total EMF εε around the loop is then:
ε = 2Blv sin θ (22.1.6)
This expression is valid, but it does not give EMF as a function of time. To find the time dependence of EMF, we assume the coil rotates at a constant angular
velocity ω. The angle θ is related to angular velocity by θ = ωt , so that:
ε = 2Blv sin ωt (22.1.7)
Now, linear velocity v is related to angular velocity by v = rω . Here r = w/2 , so that v = (w/2)ω , and:
w
ε = 2Bl ω sin ωt = (lw)Bω sin ωt (22.1.8)
2
Noting that the area of the loop is A = ℓw , and allowing for N loops, we find that:
ε = NABw sin ωt is the EMF induced in a generator coil of N turns and area A rotating at a constant angular velocity in a uniform magnetic field B.
Generators illustrated in this Atom look very much like the motors illustrated previously. This is not coincidental. In fact, a motor becomes a generator when its
shaft rotates.
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Electric Motors
learning objectives
Explain how force is generated into electric motors
An electric motor is a device that converts electrical energy into mechanical energy.
The basic principles of operation for a motor are the same as those for a generator, except that a motor converts electrical energy into mechanical energy
(motion). (Read our atom on electric generators first. ) Most electric motors use the interaction of magnetic fields and current-carrying conductors to generate
force. Electric motors are found in applications as diverse as industrial fans, blowers and pumps, machine tools, household appliances, power tools, and disk
drives.
Lorentz Force
If you were to place a moving charged particle in a magnetic field, it would experience a force called the Lorentz force:
F = q×v×B (22.1.9)
Right-Hand Rule: Right-hand rule showing the direction of the Lorentz force
where v is the speed of the moving charge, q is the charge, and B is the magnetic field. Current in a conductor consists of moving charges. Therefore, a current-
carrying coil in a magnetic field will also feel the Lorentz force. For a straight current carrying wire that is not moving, the Lorentz force is:
F = I ×L×B (22.1.10)
where F is the force (in newtons, N), I is the current in the wire (in amperes, A), L is the length of the wire that is in the magnetic field (in m), and B is the
magnetic field strength (in teslas, T). The direction of the Lorentz force is perpendicular to both the direction of the flow of current and the magnetic field and
can be found using the right-hand rule, shown in. Using your right hand, point your thumb in the direction of the current, and point your first finger in the
direction of the magnetic field. Your third finger will now be pointing in the direction of the force.
Torque: The force on opposite sides of the coil will be in opposite directions because the charges are moving in opposite directions. This means the coil will
rotate.
Mechanics of a Motor
Both motors and generators can be explained in terms of a coil that rotates in a magnetic field. In a generator the coil is attached to an external circuit that is
then turned. This results in a changing flux, which induces an electromagnetic field. In a motor, a current-carrying coil in a magnetic field experiences a force
on both sides of the coil, which creates a twisting force (called a torque) that makes it turn. Any coil carrying current can feel a force in a magnetic field. This
force is the Lorentz force on the moving charges in the conductor. The force on opposite sides of the coil will be in opposite directions because the charges are
moving in opposite directions. This means the coil will rotate.
Inductance
Inductance is the property of a device that tells how effectively it induces an emf in another device or on itself.
learning objectives
Describe properties of an inductor, distinguishing mutual inductance and self-inductance
Induction is the process in which an emf is induced by changing magnetic flux. Transformers, for example, are designed to be particularly effective at inducing
a desired voltage and current with very little loss of energy to other forms (see our Atom on “Transformers. “) Is there a useful physical quantity related to how
“effective” a given device is? The answer is yes, and that physical quantity is called inductance.
Mutual Inductance
Mutual inductance is the effect of Faraday’s law of induction for one device upon another, such as the primary coil in transmitting energy to the secondary in a
transformer. See, where simple coils induce emfs in one another.
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Mutual Inductance in Coils: These coils can induce emfs in one another like an inefficient transformer. Their mutual inductance M indicates the effectiveness
of the coupling between them. Here a change in current in coil 1 is seen to induce an emf in coil 2. (Note that “E2 induced” represents the induced emf in coil
2. )
In the many cases where the geometry of the devices is fixed, flux is changed by varying current. We therefore concentrate on the rate of change of current, ΔI/
Δt, as the cause of induction. A change in the current I1 in one device, coil 1, induces an EMF2 in the other. We express this in equation form as
ΔI1
EMF 2 = −M (22.1.11)
Δt
where M is defined to be the mutual inductance between the two devices. The minus sign is an expression of Lenz’s law. The larger the mutual inductance M,
the more effective the coupling.
Nature is symmetric here. If we change the current I2 in coil 2, we induce an emf1 in coil 1, which is given by
ΔI2
EMF 1 = −M (22.1.12)
Δt
where M is the same as for the reverse process. Transformers run backward with the same effectiveness, or mutual inductance M.
Self-Inductance
Self-inductance, the effect of Faraday’s law of induction of a device on itself, also exists. When, for example, current through a coil is increased, the magnetic
field and flux also increase, inducing a counter emf, as required by Lenz’s law. Conversely, if the current is decreased, an emf is induced that opposes the
decrease. Most devices have a fixed geometry, and so the change in flux is due entirely to the change in current ΔI through the device. The induced emf is
related to the physical geometry of the device and the rate of change of current. It is given by
ΔI
EMF = −L (22.1.13)
Δt
where L is the self-inductance of the device. A device that exhibits significant self-inductance is called an inductor. Again, the minus sign is an expression of
Lenz’s law, indicating that emf opposes the change in current.
learning objectives
Formulate two views that are applied to calculate the electromotive force
An electromotive force (EMF) induced by motion relative to a magnetic field B is called a motional EMF. You might have noticed that motional EMF is very
similar to the induced EMF caused by a changing magnetic field. In this Atom we see that they are indeed the same phenomenon, shown in different frame of
reference.
Motional EMF
In the case where a conductor loop is moving into magnet shown in (a), magnetic force on a moving charge in the loop is given by evBevB (Lorentz force, e:
electron charge).
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Conductor Loop Moving Into a Magnet: (a) Motional EMF. The current loop is moving into a stationary magnet. The direction of the magnetic field is into
the screen. (b) Induced EMF. Current loop is stationary, and the magnet is moving.
Due to the force, electrons will keep building up on one side (bottom end in the figure) until enough of an electric field opposing the motion of electrons is
established across the rod, which is eE. Equating the two forces, we get E = vB .
Therefore, the motional EMF over the length L of the side of the loop is given by εmotion = vB × L (Eq. 1), where L is the length of the object moving at
speed v relative to the magnet.
Induced EMF
Since the rate of change of the magnetic flux passing through the loop is B
dA
dt
(A: area of the loop that magnetic field pass through), the induced EMF
εinduced= BLv (Eq. 2).
learning objectives
Apply the law of conservation of energy to describe the production motional electromotive force with mechanical work
We learned about motional EMF previously (see our Atom on “Motional EMF”). For the simple setup shown below, motional EMF (ε)(ε) produced by a
moving conductor (in a uniform field) is given as follows:
ε = Blv (22.1.14)
where B is the magnetic field, l is the length of the conducting rod, and v is the (constant) speed of its motion. (B, l, and v are all perpendicular to each other as
shown in the image below.)
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Motional EMF: (a) A motional emf=Bℓv is induced between the rails when this rod moves to the right in the uniform magnetic field. The magnetic field B is
into the page, perpendicular to the moving rod and rails and, hence, to the area enclosed by them. (b) Lenz’s law gives the directions of the induced field and
current, and the polarity of the induced emf. Since the flux is increasing, the induced field is in the opposite direction, or out of the page. Right hand rule gives
the current direction shown, and the polarity of the rod will drive such a current.
Conservation of Energy
In this atom, we will consider the system from the energy perspective. As the rod moves and carries current i, it will feel the Lorentz force
FL = iBL (22.1.15)
To keep the rod moving at a constant speed v, we must constantly apply an external force Fext(equal to magnitude of FL and opposite in its direction) to the rod
along its motion. Since the rod is moving at v, the power P delivered by the external force would be:
P = Fext v = (iBL) × v = iε (22.1.16)
In the final step, we used the first equation we talked about. Note that this is exactly the power dissipated in the loop (= current × voltage ). Therefore, we
conclude that the mechanical work done by an external force to keep the rod moving at a constant speed is converted to heat energy in the loop. More
generally, mechanical work done by an external force to produce motional EMF is converted to heat energy. Energy is conserved in the process.
Lenz’ Law
We learned in the Atom “Faraday’s Law of Induction and Lenz’ Law” that Lenz’ law is a manifestation of the conservation of energy. As we see in the
example in this Atom, Lenz’ law guarantees that the motion of the rod is opposed because of nature’s tendency to oppose a change in magnetic field. If the
induced EMF were in the same direction as the change in flux, there would be a positive feedback causing the rod to fly away from the slightest perturbation.
2μ
.
learning objectives
Express the energy density of a magnetic field in a form of equation
Energy is needed to generate a magnetic field both to work against the electric field that a changing magnetic field creates and to change the magnetization of
any material within the magnetic field. For non-dispersive materials this same energy is released when the magnetic field is destroyed. Therefore, this energy
can be modeled as being “stored” in the magnetic field.
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Magnetic Field Created By A Solenoid: Magnetic field created by a solenoid (cross-sectional view) described using field lines. Energy is “stored” in the
magnetic field.
Energy density is the amount of energy stored in a given system or region of space per unit volume. If there are no magnetic materials around, μ can be
replaced by μ0. The above equation cannot be used for nonlinear materials, though; a more general expression (given below) must be used.
In general, the incremental amount of work per unit volume δW needed to cause a small change of magnetic field δB is:
δW = H ⋅ δB (22.1.18)
Once the relationship between H and B is known this equation is used to determine the work needed to reach a given magnetic state. For hysteretic materials
such as ferromagnets and superconductors, the work needed also depends on how the magnetic field is created. For linear non-dispersive materials, though, the
general equation leads directly to the simpler energy density equation given above.
Proof: Power that should be supplied to an inductor with inductance L to run current I through it it given as
dI
P = VI = L ×I (22.1.20)
dt
Therefore
T I
′ ′
1 2
Estored = ∫ P(t)dt = ∫ LI dI = LI (22.1.21)
0 0
2
Transformers
Transformers transform voltages from one value to another; its function is governed by the transformer equation.
learning objectives
Apply the transformer equation to compare the secondary and primary voltages
Transformers change voltages from one value to another. For example, devices such as cell phones, laptops, video games, power tools and small appliances
have a transformer (built into their plug-in unit) that changes 120 V into the proper voltage for the device. Transformers are also used at several points in power
distribution systems, as shown in. Power is sent long distances at high voltages, as less current is required for a given amount of power (this means less line
loss). Because high voltages pose greater hazards, transformers are employed to produce lower voltage at the user’s location.
Transformer Setup: Transformers change voltages at several points in a power distribution system. Electric power is usually generated at greater than 10 kV,
and transmitted long distances at voltages over 200 kV—sometimes as great as 700 kV—to limit energy losses. Local power distribution to neighborhoods or
industries goes through a substation and is sent short distances at voltages ranging from 5 to 13 kV. This is reduced to 120, 240, or 480 V for safety at the
individual user site.
The type of transformer considered here is based on Faraday’s law of induction, and is very similar in construction to the apparatus Faraday used to
demonstrate that magnetic fields can create currents (illustrated in ). The two coils are called the primary and secondary coils. In normal use, the input voltage
is placed on the primary, and the secondary produces the transformed output voltage. Not only does the iron core trap the magnetic field created by the primary
coil, its magnetization increases the field strength. Since the input voltage is AC, a time-varying magnetic flux is sent to the secondary, inducing its AC output
voltage.
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Simple Transformer: A typical construction of a simple transformer has two coils wound on a ferromagnetic core that is laminated to minimize eddy currents.
The magnetic field created by the primary is mostly confined to and increased by the core, which transmits it to the secondary coil. Any change in current in
the primary induces a current in the secondary.The figure shows a simple transformer with two coils wound on either sides of a laminated ferromagnetic core.
The set of coil on left side of the core is marked as the primary and there number is given as N p. The voltage across the primary is given by V p. The set of
coil on right side of the core is marked as the secondary and there number is represented as N s. The voltage across the secondary is given by V s. A symbol of
the transformer is also shown below the diagram. It consists of two inductor coils separated by two equal parallel lines representing the core.
Transformer Equation
For the simple transformer shown in, the output voltage Vsdepends almost entirely on the input voltage Vp and the ratio of the number of loops in the primary
and secondary coils. Faraday’s law of induction for the secondary coil gives its induced output voltage Vs as:
ΔΦ
Vs = −Ns (22.1.22)
Δt
where Ns is the number of loops in the secondary coil and Δ/Δt is the rate of change of magnetic flux. Note that the output voltage equals the induced EMF
(Vs=EMFs), provided coil resistance is small. The cross-sectional area of the coils is the same on either side, as is the magnetic field strength, so /Δt is the same
on either side. The input primary voltage Vp is also related to changing flux by:
ΔΦ
Vp = −N[ (22.1.23)
Δt
Taking the ratio of these last two equations yields a useful relationship:
Vs Ns
= (22.1.24)
Vp Np
This is known as the transformer equation, which simply states that the ratio of the secondary to primary voltages in a transformer equals the ratio of the
number of loops in their coils. The output voltage of a transformer can be less than, greater than or equal to the input voltage, depending on the ratio of the
number of loops in their coils. Some transformers even provide a variable output by allowing connection to be made at different points on the secondary coil. A
step-up transformer is one that increases voltage, whereas a step-down transformer decreases voltage.
Assuming, as we have, that resistance is negligible, the electrical power output of a transformer equals its input. Equating the power input and output,
Pp = Ip Vp = Is Vs = Ps (22.1.25)
Key Points
It is a change in the magnetic field flux that results in an electromotive force (or voltage).
The magnetic flux (often denoted Φ or ΦB) through a surface is the component of the magnetic field passing through that surface.
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In the most general form, magnetic flux is defined as Φ = ∬ B ⋅ dA . It is the integral (sum) of all of the magnetic field passing through infinitesimal
B
A
∂t
Faraday’s law of induction is one of the four equations in Maxwell’s equations, governing all electromagnetic phenomena.
An electric generator rotates a coil in a magnetic field, inducing an EMF given as a function of time by ε = NABw sin ωt .
Generators supply almost all of the power for the electric power grids which provide most of the world’s electric power.
A motor becomes a generator when its shaft rotates.
Most electric motors use the interaction of magnetic fields and current -carrying conductors to generate force.
Current in a conductor consists of moving charges. Therefore, a current-carrying coil in a magnetic field will also feel the Lorentz force.
In a motor, a current-carrying coil in a magnetic field experiences a force on both sides of the coil, which creates a twisting force (called a torque) that
makes it turn.
Mutual inductance is the effect of two devices in inducing emfs in each other. A change in current ΔI1/Δt in one induces an emf emf2 in the seccond:
EMF2= −M ΔI1/Δt, where M is defined to be the mutual inductance between the two devices.
Self-inductance is the effect of the device inducing emf in itself.
A device that exhibits significant self-inductance is called an inductor, and the EMF induced in it by a change in current through it is EMF = −L . ΔI
Δt
Motional and induced EMF are the same phenomenon, just observed in different reference frames. Equivalence of the two phenomena is what triggered
Einstein to work on special relativity.
The EMF produced due to the relative motion of the loop and magnet is given as ε = vB × L (Eq. 1), where L is the length of the object moving at
motion
2
2
Transformers are often used at several points in the power distribution systems and also in many household power adapters.
T ransf ormerequationstatesthattheratioof thesecondarytoprimaryvoltagesinatransf ormerequalstheratioof thenumberof loopsintheircoils
Vs Ns
: \( =
Vp Np
.
Assuming, as we have, that resistance is negligible, the electrical power output of a transformer equals its input. This leads us to another useful eqaution:
Is Np
Ip
=
Ns
. If voltage increases, current decreases. Conversely, if voltage decreases, current increases.
Key Terms
vector area: A vector whose magnitude is the area under consideration, and whose direction is perpendicular to the surface area.
galvanometer: An analog measuring device, denoted by G, that measures current flow using a needle deflection caused by a magnetic field force acting
upon a current-carrying wire.
electromotive force: (EMF)—The voltage generated by a battery or by the magnetic force according to Faraday’s Law. It is measured in units of volts, not
newtons, and thus, is not actually a force.
solenoid: A coil of wire that acts as a magnet when an electric current flows through it.
flux: The rate of transfer of energy (or another physical quantity) through a given surface, specifically electric flux or magnetic flux.
magnetic flux: A measure of the strength of a magnetic field in a given area.
induction: The generation of an electric current by a varying magnetic field.
Faraday’s law of induction: A basic law of electromagnetism that predicts how a magnetic field will interact with an electric circuit to produce an
electromotive force (EMF).
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Maxwell’s equations: A set of equations describing how electric and magnetic fields are generated and altered by each other and by charges and currents.
Stokes’ theorem: a statement about the integration of differential forms on manifolds, which both simplifies and generalizes several theorems from vector
calculus.
turbine: Any of various rotary machines that use the kinetic energy of a continuous stream of fluid (a liquid or a gas) to turn a shaft.
Lorentz force: The force exerted on a charged particle in an electromagnetic field.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
transformer: A static device that transfers electric energy from one circuit to another by magnetic coupling. Their main use is to transfer energy between
different voltage levels, which allows choosing most appropriate voltage for power generation, transmission and distribution separately.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that the speed of light is
constant in all frames of reference.
magnetic field: A condition in the space around a magnet or electric current in which there is a detectable magnetic force, and where two magnetic poles
are present.
frame of reference: A coordinate system or set of axes within which to measure the position, orientation, and other properties of objects in it.
motional EMF: An EMF (electromotive force) induced by motion relative to a magnetic field.
permeability: A quantitative measure of the degree of magnetization of a material in the presence of an applied magnetic field (measured in newtons per
ampere squared in SI units).
inductor: A passive device that introduces inductance into an electrical circuit.
ferromagnet: Materials that show a permanent magnetic property.
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Solenoid. Provided by: Wikipedia. Located at: en.Wikipedia.org/wiki/Solenoid. License: CC BY: Attribution
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22.1: Magnetic Flux, Induction, and Faraday’s Law is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
22.1.18 https://phys.libretexts.org/@go/page/15654
22.2: AC Circuits
learning objectives
Describe properties of an inductor
Inductance
OVERVIEW
Induction is the process in which an emf is induced by changing magnetic flux. Specifically in the case of electronics, inductance is
the property of a conductor by which a change in current in the conductor creates a voltage in both the conductor itself (self-
inductance) and any nearby conductors (mutual inductance). This effect derives from two fundamental observations of physics:
First, that a steady current creates a steady magnetic field and second, that a time-varying magnetic field induces a voltage in a
nearby conductor (Faraday’s law of induction). From Lenz’s law, a changing electric current through a circuit that has inductance
induces a proportional voltage which opposes the change in current (if this wasn’t true one can easily see how energy could not be
conserved, with a changing current reinforcing the change in a positive feedback loop).
MUTUAL INDUCTANCE
Mutual inductance is the effect of Faraday’s law of induction for one device upon another, such as the primary coil in transmitting
energy to the secondary in a transformer. See, where simple coils induce emfs in one another.
Mutual Inductance in Coils: These coils can induce emfs in one another like an inefficient transformer. Their mutual inductance
M indicates the effectiveness of the coupling between them. Here a change in current in coil 1 is seen to induce an emf in coil 2.
(Note that “E2 induced” represents the induced emf in coil 2. )
In the many cases where the geometry of the devices is fixed, flux is changed by varying current. We therefore concentrate on the
rate of change of current, ΔI/Δt, as the cause of induction. A change in the current I1 in one device, coil 1 in the figure, induces an
emf2 in the other. We express this in equation form as
ΔI1
emf2 = −M (22.2.1)
Δt
where M is defined to be the mutual inductance between the two devices. The minus sign is an expression of Lenz’s law. The larger
the mutual inductance M, the more effective the coupling. Units for M are (V⋅s)/A=Ωs, which is named a henry (H), after Joseph
Henry (discovered of self-inductance). That is, 1 H=1Ωs.
Nature is symmetric here. If we change the current I2 in coil 2, we induce an emf1 in coil 1, which is given by
ΔI2
emf1 = −M (22.2.2)
Δt
where M is the same as for the reverse process. Transformers run backward with the same effectiveness, or mutual inductance M.
A large mutual inductance M may or may not be desirable. We want a transformer to have a large mutual inductance. But an
appliance, such as an electric clothes dryer, can induce a dangerous emf on its case if the mutual inductance between its coils and
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the case is large. One way to reduce mutual inductance M is to counterwind coils to cancel the magnetic field produced. (See ).
Counterwinding: The heating coils of an electric clothes dryer can be counter-wound so that their magnetic fields cancel one
another, greatly reducing the mutual inductance with the case of the dryer.
SELF-INDUCTANCE
Self-inductance, the effect of Faraday’s law of induction of a device on itself, also exists. When, for example, current through a coil
is increased, the magnetic field and flux also increase, inducing a counter emf, as required by Lenz’s law. Conversely, if the current
is decreased, an emf is induced that opposes the decrease. Most devices have a fixed geometry, and so the change in flux is due
entirely to the change in current ΔI through the device. The induced emf is related to the physical geometry of the device and the
rate of change of current. It is given by
ΔI
emf = −L (22.2.3)
Δt
where L is the self-inductance of the device. A device that exhibits significant self-inductance is called an inductor, and given the
symbol in.
Inductor Symbol
The minus sign is an expression of Lenz’s law, indicating that emf opposes the change in current. Units of self-inductance are
henries (H) just as for mutual inductance. The larger the self-inductance L of a device, the greater its opposition to any change in
current through it. For example, a large coil with many turns and an iron core has a large L and will not allow current to change
quickly. To avoid this effect, a small L must be achieved, such as by counterwinding coils as in.
SOLENOIDS
It is possible to calculate L for an inductor given its geometry (size and shape) and knowing the magnetic field that it produces.
This is difficult in most cases, because of the complexity of the field created. The inductance L is usually a given quantity. One
exception is the solenoid, because it has a very uniform field inside, a nearly zero field outside, and a simple shape. The self-
inductance of a solenoid of cross-sectional area A and length ℓ is
2
μ0 N A
L = (solenoid). (22.2.4)
l
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It is instructive to derive this equation, but this is left as an exercise to the reader. (Hint: start by noting that the induced emf is
given by Faraday’s law of induction as emf=−N(Δ/Δt) and, by the definition of self-inductance is given as as emf=−L(ΔI//Δt) and
equate these two expressions). Note that the inductance depends only on the physical characteristics of the solenoid, consistent with
its definition.
RL Circuits
An RL circuit consists of an inductor and a resistor, in series or parallel with each other, with current driven by a voltage source.
learning objectives
Describe current-voltage relationship in the RL circuit and calculate energy that can be stored in an inductor
RL Circuits
A resistor-inductor circuit (RL circuit) consists of a resistor and an inductor (either in series or in parallel ) driven by a voltage
source.
Review
Recall that induction is the process in which an emf is induced by changing magnetic flux. Mutual inductance is the effect of
Faraday’s law of induction for one device upon another, while self-inductance is the the effect of Faraday’s law of induction of a
device on itself. An inductor is a device or circuit component that exhibits self-inductance.
Energy of an Inductor
We know from Lenz’s law that inductors oppose changes in current. We can think of this situation in terms of energy. Energy is
stored in a magnetic field. It takes time to build up energy, and it also takes time to deplete energy; hence, there is an opposition to
rapid change. In an inductor, the magnetic field is directly proportional to current and to the inductance of the device. It can be
shown that the energy stored in an inductor Eind is given by:
1
2
Eind = LI (22.2.5)
2
Inductors in Circuits
We know that the current through an inductor L cannot be turned on or off instantaneously. The change in current changes the
magnetic flux, inducing an emf opposing the change (Lenz’s law). How long does the opposition last? Current will flow and can be
turned off, but how long does it take? The following figure shows a switching circuit that can be used to examine current through
an inductor as a function of time.
Current in an RL Circuit: (a) An RL circuit with a switch to turn current on and off. When in position 1, the battery, resistor, and
inductor are in series and a current is established. In position 2, the battery is removed and the current eventually stops because of
energy loss in the resistor. (b) A graph of current growth versus time when the switch is moved to position 1. (c) A graph of current
decay when the switch is moved to position 2.
When the switch is first moved to position 1 (at t=0), the current is zero and it eventually rises to I0=V/R, where R is the total
resistance of the circuitand V is the battery ‘s voltage. The opposition of the inductor L is greatest at the beginning, because the
change in current is greatest at that time. The opposition it poses is in the form of an induced emf, which decreases to zero as the
current approaches its final value. This is the hallmark of an exponential behavior, and it can be shown (with calculus) that
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−t
⎛ ⎞
I = I0 ⎜1 − e τ ⎟ (22.2.6)
⎝ ⎠
is the current in an RL circuit when switched on. (Note the similarity to the exponential behavior of the voltage on a charging
capacitor.) The initial current is zero and approaches I0=V/R with a characteristic time constant for an RL circuit, given by:
L
τ = (22.2.7)
R
where ττ has units of seconds, since 1H=1Ω⋅s1H=1Ω⋅s. In the first period of time ττ, the current rises from zero to 0.632I0, since
I=I0(1−e−1)=I0(1−0.368)=0.632I0. The current will be 0.632 of the remainder in the next time. A well-known property of the
exponential function is that the final value is never exactly reached, but 0.632 of the remainder to that value is achieved in every
characteristic time ττ. In just a few multiples of the time ττ, the final value is very nearly achieved (see part (b) of above figure).
The characteristic time ττ depends on only two factors, the inductance L and the resistance R. The greater the inductance L, the
greater it is, which makes sense since a large inductance is very effective in opposing change. The smaller the resistance R, the
greater ττ is. Again this makes sense, since a small resistance means a large final current and a greater change to get there. In both
cases (large Land small R) more energy is stored in the inductor and more time is required to get it in and out.
When the switch in (a) is moved to position 2 and cuts the battery out of the circuit, the current drops because of energy dissipation
by the resistor. However, this is also not instantaneous, since the inductor opposes the decrease in current by inducing an emf in the
same direction as the battery that drove the current. Furthermore, there is a certain amount of energy, (1/2)LI02, stored in the
inductor, and it is dissipated at a finite rate. As the current approaches zero, the rate of decrease slows, since the energy dissipation
rate is I2R. Once again the behavior is exponential, and I is found to be
−t
I = I0 e τ
(22.2.8)
In (c), in the first period of time τ=L/Rτ=L/R after the switch is closed, the current falls to 0.368 of its initial value, since
I=I0e−1=0.368I0. In each successive time ττ, the current falls to 0.368 of the preceding value, and in a few multiples of ττ, the
current becomes very close to zero.
learning objectives
Distinguish behavior of RLC series circuits as large and small frequencies
In previous Atoms we learned how an RLC series circuit, as shown in, responds to an AC voltage source. By combining Ohm’s law
(Irms=Vrms/Z; Irms and Vrms are rms current and voltage) and the expression for impedance Z, from:
Series RLC Circuit: A series RLC circuit: a resistor, inductor and capacitor (from left).
−−−−−−−−−−−−−−
2
1
2
Z = √R + (XL − XC ) ( XL = 2πν L, XC = ) (22.2.9)
2πν C
Vr m s
we arrived at: I
rms =
2
.
2
√R +( XL −XC )
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From the equation, we studied resonance conditions for the circuit. We also learned the phase relationships among the voltages
across resistor, capacitor and inductor: when a sinusoidal voltage is applied, the current lags the voltage by a 90º phase in a circuit
with an inductor, while the current leads the voltage by 90∘ in a circuit with a capacitor. Now, we will examine the system’s
response at limits of large and small frequencies.
At Large Frequencies
At large enough frequencies (ν ≫
1
√2πLC
), XL is much greater than XC. If the frequency is high enough that XL is much larger
than R as well, the impedance Z is dominated by the inductive term. When Z ≈ X , the circuit is almost equivalent to an AC
L
circuit with just an inductor. Therefore, the rms current will be Vrms/XL, and the current lags the voltage by almost 90∘. This
response makes sense because, at high frequencies, Lenz’s law suggests that the impedance due to the inductor will be large.
At Small Frequencies
The impedance Z at small frequencies (ν ≪
1
) is dominated by the capacitive term, assuming that the frequency is high
√2πLC
enough so that XC is much larger than R. When Z ≈ X , the circuit is almost equivalent to an AC circuit with just a capacitor.
C
Therefore, the rms current will be given as Vrms/XC, and the current leads the voltage by almost 90∘.
Resistors in AC Circuits
In a circuit with a resistor and an AC power source, Ohm’s law still applies (V = IR).
learning objectives
Apply Ohm’s law to determine current and voltage in an AC circuit
Direct current (DC) is the flow of electric charge in only one direction. It is the steady state of a constant-voltage circuit. Most well
known applications, however, use a time-varying voltage source. Alternating current (AC) is the flow of electric charge that
periodically reverses direction. If the source varies periodically, particularly sinusoidally, the circuit is known as an alternating-
current circuit. Examples include the commercial and residential power that serves so many of our needs. shows graphs of voltage
and current versus time for typical DC and AC power. The AC voltages and frequencies commonly used in homes and businesses
vary around the world.
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Sinusoidal Voltage and Current: (a) DC voltage and current are constant in time, once the current is established. (b) A graph of
voltage and current versus time for 60-Hz AC power. The voltage and current are sinusoidal and are in phase for a simple
resistance circuit. The frequencies and peak voltages of AC sources differ greatly.
We have studied Ohm’s law:
V
I = (22.2.10)
R
where I is the current, V is the voltage, and R is the resistance of the circuit. Ohm’s law applies to AC circuits as well as to DC
circuits. Therefore, with an AC voltage given by:
V = V0 sin(2πν t) (22.2.11)
where V0 is the peak voltage and νν is the frequency in hertz, the current in the circuit is given as:
V0
I = sin(2πν t) (22.2.12)
R
In this example, in which we have a resistor and the voltage source in the circuit, the voltage and current are said to be in phase, as
seen in (b). Current in the resistor alternates back and forth without any phase difference, just like the driving voltage.
Consider a perfect resistor that brightens and dims 120 times per second as the current repeatedly goes through zero. (A 120-Hz
flicker is too rapid for your eyes to detect. ) The fact that the light output fluctuates means that the power is fluctuating. Since the
power supplied is P = IV, if we use the above expressions for I and V, we see that the time dependence of power is:
2
V0
P = ⋅ sin(2πν t) (22.2.13)
R
To find the average power consumed by this circuit, we need to take the time average of the function. Since:
π
1 2
1
∫ sin (x)dx = (22.2.14)
π 0
2
we see that:
2
V0
Pavg = (22.2.15)
2R
learning objectives
Explain the benefits of using a phasor representation
In the previous Atom on “Resistors in AC Circuits”, we introduced an AC power source and studied how resistors behave in AC
circuits. There, we used the Ohm’s law (V=IR) to derive the relationship between voltage and current in AC circuits. In this and
following Atoms, we will generalize the Ohm’s law so that we can use it even when we have capacitors and inductors in the circuit.
To get there, we will first introduce a very general, pictorial way of representing a sinusoidal wave, using phasor.
22.2.6 https://phys.libretexts.org/@go/page/15656
Capacitors in AC Circuits with Phasors | Doc Physics
Phasor
The key idea in the phasor representation is that a complex, time-varying signal may be represented as the product of a complex
number (that is independent of time) and a complex signal (that is dependent on time). Phasors separate the dependencies on A
(amplitude), νν ( frequency ), and θ ( phase ) into three independent factors. This can be particularly useful because the frequency
factor (which includes the time-dependence of the sinusoid) is often common to all the components of a linear combination of
sinusoids. In those situations, phasors allow this common feature to be factored out, leaving just the A and θ features. For example,
we can represent A ⋅ cos(2πν t + θ) simply as a complex constant, Ae . Since phasors are represented by a magnitude (or
iθ
modulus) and an angle, it is pictorially represented by a rotating arrow (or a vector) in x-y plane.
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Fig 3: A phasor can be seen as a vector rotating about the origin in a complex plane. The cosine function is the projection of the
vector onto the real axis. Its amplitude is the modulus of the vector, and its argument is the total phase \omega t+\theta. The phase
constant \theta represents the angle that the vector forms with the real axis at t = 0.
Capacitors in AC circuits
Earlier in a previous Atom, we studied how the voltage and the current varied with time. If the AC supply is connected to a resistor,
then the current and voltage will be proportional to each other. This means that the current and voltage will “peak” at the same
time. We say that the current and voltage are in phase.
When a capacitor is connected to an alternating voltage, the maximum voltage is proportional to the maximum current, but the
maximum voltage does not occur at the same time as the maximum current. The current has its maximum (it peaks) one quarter of
a cycle before the voltage peaks. Engineers say that the “current leads the voltage by 90∘“. This is shown in.
Fig 2: The current peaks (has its maximum) one quarter of a wave before the voltage when a capacitor is connected to an
alternating voltage.
For a circuit with a capacitor, the instantaneous value of V/I is not constant. However, the value of Vmax/Imax is useful, and is called
the capacitive reactance (XC) of the component. Because it is still a voltage divided by a current (like resistance ), its unit is the
ohm. The value of XC (C standing for capacitor) depends on its capacitance (C) and the frequency (f) of the alternating current.
X C =
1
2πν C
.
The capacitor is affecting the current, having the ability to stop it altogether when fully charged. Since an AC voltage is applied,
there is an rms current, but it is limited by the capacitor. This is considered to be an effective resistance of the capacitor to AC, and
Vr m s
so the rms current Irms in the circuit containing only a capacitor C is given by another version of Ohm’s law to be I rms =
XC
,
where Vrms is the rms voltage. Note that XC replaces R in the DC version of the Ohm’s law.
Phase representation
Since the voltage across a capacitor lags the current, the phasor representing the current and voltage would be give as in. In the
diagram, the arrows rotate in counter-clockwise direction at a frequency νν. (Therefore, current leads voltage. ) In the following
Atoms, we will see how these phasors can be used to analyze RC, RL, LC, and RLC circuits.
22.2.8 https://phys.libretexts.org/@go/page/15656
Fig 4: Phasor diagram for an AC circuit with a capacitor
learning objectives
Explain why the voltage across an inductor “leads” the current in an AC circuit with an inductor
Suppose an inductor is connected directly to an AC voltage source, as shown in. It is reasonable to assume negligible resistance
because in practice we can make the resistance of an inductor so small that it has a negligible effect on the circuit. The graph shows
voltage and current as functions of time. (b) starts with voltage at a maximum. Note that the current starts at zero, then rises to its
peak after the voltage driving it (as seen in the preceding section when DC voltage was switched on).
AC Voltage Source in Series with an Inductor: (a) An AC voltage source in series with an inductor having negligible resistance.
(b) Graph of current and voltage across the inductor as functions of time.
When the voltage becomes negative at point a, the current begins to decrease; it becomes zero at point b, where voltage is its most
negative. The current then becomes negative, again following the voltage. The voltage becomes positive at point c where it begins
to make the current less negative. At point d, the current goes through zero just as the voltage reaches its positive peak to start
another cycle. Hence, when a sinusoidal voltage is applied to an inductor, the voltage leads the current by one-fourth of a cycle, or
by a 90º phase angle.
Current lags behind voltage, since inductors oppose change in current. Changing current induces an emf. This is considered an
effective resistance of the inductor to AC. The rms current Irms through an inductor L is given by a version of Ohm’s law:
Vr m s
I
rms =
XL
where Vrms is the rms voltage across the inductor and X = 2πν L with ν the frequency of the AC voltage source in
L
hertz. XL is called the inductive reactance. Because the inductor reacts to impede the current, XLhas units of ohms (1 H=1 Ωs, so
that frequency times inductance has units of (cycles/s)(Ωs)=Ω), consistent with its role as an effective resistance.
Phasor Representation
The voltage across an inductor “leads” the current because of the Lenz’s law. Therefore, the phasor representing the current and
voltage would be given as in. Again, the phasors are vectors rotating in counter-clockwise direction at a frequency ν (you can see
that the voltage leads the current). Subsequent Atoms will discuss how these phasors can be used to analyze RC, RL, LC, and RLC
circuits.
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Phasor Diagram: Phasor diagram for an AC circuit with an inductor.
22.2.10 https://phys.libretexts.org/@go/page/15656
learning objectives
Compare resonance characteristics of higher- and lower-resistance circuits
Resonance is the tendency of a system to oscillate with greater amplitude at some frequencies than at others. Frequencies at which
the response amplitude is a relative maximum are known as the system’s resonance frequencies. To study the resonance in an RLC
circuit, as illustrated below, we can see how the circuit behaves as a function of the frequency of the driving voltage source.
RLC Series Circuit: An RLC series circuit with an AC voltage source. f is the frequency of the source.
Combining Ohm ‘s law, Irms=Vrms/Z, and the expression for impedance Z from
−−−−−−−−−−−−−−
Z = √R
2
+ (XL − XC )
2
gives
Vrms
Irms = −−−−−−−−−−−−−− (22.2.16)
2 2
√R + (XL − XC )
where Irms and Vrms are rms current and voltage, respectively. The reactances vary with frequency νν, with XL large at high
frequencies and XC large at low frequencies given as:
1
XL = 2πν L, XC = (22.2.17)
2πν C
At some intermediate frequency ν , the reactances will be equal and cancel, giving Z=R —this is a minimum value for impedance,
0
and a maximum value for Irms results. We can get an expression for ν0ν0 by taking XL=XC. Substituting the definitions of XL and
XC yields:
1
ν0 = −
−− (22.2.18)
2π √LC
ν0 is the resonant frequency of an RLC series circuit. This is also the natural frequency at which the circuit would oscillate if not
driven by the voltage source. At ν , the effects of the inductor and capacitor cancel, so that Z=R, and Irms is a maximum.
0
Resonance in AC circuits is analogous to mechanical resonance, where resonance is defined as a forced oscillation (in this case,
forced by the voltage source) at the natural frequency of the system.
The receiver in a radio is an RLC circuit that oscillates best at its ν . A variable capacitor is often used to adjust the resonance
0
frequency to receive a desired frequency and to reject others. is a graph of current as a function of frequency, illustrating a resonant
peak in Irmsat ν = f . The two curves are for two different circuits, which differ only in the amount of resistance in them. The
0 0
peak is lower and broader for the higher-resistance circuit. Thus higher-resistance circuits do not resonate as strongly, nor would
they be as selective in, for example, a radio receiver.
22.2.11 https://phys.libretexts.org/@go/page/15656
Current vs. Frequency: A graph of current versus frequency for two RLC series circuits differing only in the amount of resistance.
Both have a resonance at f0, but that for the higher resistance is lower and broader. The driving AC voltage source has a fixed
amplitude V0.
Power
Power delivered to an RLC series AC circuit is dissipated by the resistance in the circuit, and is given as Pavg = Irms Vrms cos ϕ .
Here, ϕ is called the phase angle.
learning objectives
Calculate the power delivered to an RLC-series AC circuit given the current and the voltage
If current varies with frequency in an RLC circuit, then the power delivered to it also varies with frequency. However, the average
power is not simply current times voltage, as is the case in purely resistive circuits. As seen in previous Atoms, voltage and current
are out of phase in an RLC circuit. There is a phase angle ϕ between the source voltage V and the current I, given as
cos ϕ =
R
Z
, as diagramed in.
Phasor Diagram for an RLC Series Circuit: Phasor diagram for an RLC series circuit. \phi is the phase angle, equal to the phase
difference between the voltage and current.
For example, at the resonant frequency (ν 0 =
1
2π√LC
) or in a purely resistive circuit, Z=R, so that cosϕ=1. This implies that ϕ=0º
and that voltage and current are in phase, as expected for resistors. At other frequencies, average power is less than at resonance,
because voltage and current are out of phase and Irms is lower.
The fact that source voltage and current are out of phase affects the power delivered to the circuit. It can be shown that the average
power is
Pavg = Irms Vrms cos ϕ (22.2.19)
22.2.12 https://phys.libretexts.org/@go/page/15656
(an equation derived by taking a time average of power, P(t) = I(t)V(t), over a period. I(t) and V(t) are current and voltage at time
t). Thus cosϕ is called the power factor, which can range from 0 to 1. Power factors near 1 are desirable when designing an
efficient motor, for example. At the resonant frequency, cosϕ=1.
Power delivered to an RLC series AC circuit is dissipated by the resistance alone. The inductor and capacitor have energy input and
output, but do not dissipate energy out of the circuit. Rather, they transfer energy back and forth to one another, with the resistor
dissipating the exact amount that the voltage source gives the circuit. This assumes no significant electromagnetic radiation from
the inductor and capacitor (such as radio waves).
The circuit is analogous to the wheel of a car driven over a corrugated road, as seen in. The regularly spaced bumps in the road are
analogous to the voltage source, driving the wheel up and down. The shock absorber is analogous to the resistance damping and
limiting the amplitude of the oscillation. Energy within the system goes back and forth between kinetic (analogous to maximum
current, and energy stored in an inductor) and potential energy stored in the car spring (analogous to no current, and energy stored
in the electric field of a capacitor). The amplitude of the wheels’ motion is a maximum if the bumps in the road are hit at the
resonant frequency.
Forced Damped Motion of a Wheel on a Car Spring: The forced but damped motion of the wheel on the car spring is analogous
to an RLC series AC circuit. The shock absorber damps the motion and dissipates energy, analogous to the resistance in an RLC
circuit. The mass and spring determine the resonant frequency.
Key Points
In the case of electronics, inductance is the property of a conductor by which a change in current in the conductor creates a
voltage in both the conductor itself, called self-inductance, and any nearby conductors, called mutual inductance.
From Lenz’s law, a changing electric current through a circuit that has inductance induces a proportional voltage which opposes
the change in current.
Mutual inductance is illustrated by. A change in the current I1 in one device, coil 1 in the figure, induces an emf2 in the other.
ΔI1
We express this in equation form as emf = −M 2 . M is the same for the reverse process.
Δt
Self-inductance is the effect of Faraday’s law of induction of a device on itself. The induced emf is related to the physical
geometry of the device and the rate of change of current given by emf = −L . ΔI
Δt
A device that exhibits significant self-inductance is called an inductor, and given the symbol in.
The energy stored in an inductor is E = LI . It takes time to build up stored energy in a conductor and time to deplete it.
1
2
2
When a resistor and an inductor in series are connected to a voltage source, the time-dependent current is given by
−t
I = I0 (1 − e τ
) . The final current after a long time is I .
0
The characteristic time constant is given by τ = , where R is resistance and L is inductance. This represents the time
L
necessary for the current in a circuit just closed to go from zero to 0.632 ⋅ I . 0
−t
When the voltage source is disconnected from the inductor, the current will decay according to I = I e . In the first time
0
τ
e
0
RLC circuits can be described by the (generalized) Ohm ‘s law. As for the phase, when a sinusoidal voltage is applied, the
current lags the voltage by a 90º phase in a circuit with an inductor, while the current leads the voltage by 90∘ in a circuit with a
capacitor.
22.2.13 https://phys.libretexts.org/@go/page/15656
At large enough frequencies (ν ≫
1
), the circuit is almost equivalent to an AC circuit with just an inductor. Therefore, the
√2πLC
rms current will be Vrms/XL, and the current lags the voltage by almost 90∘.
At small enough frequencies (ν ≪ 1
), the circuit is almost equivalent to an AC circuit with just a capacitor. Therefore, the
√2πLC
rms current will be given as Vrms/XC, and the current leads the voltage by almost 90∘.
V0
With an AC voltage given by: V = V sin(2πν t) the current in the circuit is given as: I = sin(2πν t) This expression comes
0
R
2R
When a capacitor is connected to an alternating voltage, the maximum voltage is proportional to the maximum current, but the
maximum voltage does not occur at the same time as the maximum current.
If the AC supply is connected to a resistor, then the current and voltage will be proportional to each other. This means that the
current and voltage will “peak” at the same time.
Vr m s
The rms current in the circuit containing only a capacitor C is given by another version of Ohm ‘s law to be I = rms , where
XC
With an inductor in an AC circuit, the voltage leads the current by one-fourth of a cycle, or by a 90º phase angle.
Vr m s
The rms current Irms through an inductor L is given by a version of Ohm ‘s law: I =
rms . XL is called the inductive
XL
Phasors are vectors rotating in counter-clockwise direction. A phasor for an inductor shows that the voltage lead the current by
a 90º phase.
Resonance condition of an RLC series circuit can be obtained by equating XL and XC, so that the two opposing phasors cancel
each other.
At resonance, the effects of the inductor and capacitor cancel, so that Z=R, and Irms is a maximum.
Higher- resistance circuits do not resonate as strongly compared to lower-resistance circuits, nor would they be as selective in,
for example, a radio receiver.
Phase angle ϕ is the phase difference between the source voltage V and the current I. See the phasor diagram in.
At the resonant frequency or in a purely resistive circuit Z=R, so that cosϕ=1. This implies that ϕ=0º and that voltage and
current are in phase.
Average power dissipated in an RLC circuit can be calculated by taking a time average of power, P(t) = I(t)V(t), over a period.
Key Terms
mutual inductance: The ratio of the voltage in a circuit to the change in current in a neighboring circuit.
self-inductance: The ratio of the voltage to the change in current in the same circuit.
inductor: A passive device that introduces inductance into an electrical circuit.
characteristic time constant: Denoted by τ , in RL circuits it is given by τ = where R is resistance and L is inductance.
L
When a switch is closed, it is the time it takes for the current to decay by a factor of 1/e.
inductor: A device or circuit component that exhibits significant self-inductance; a device which stores energy in a magnetic
field.
Lenz’s law: A law of electromagnetic induction that states that an electromotive force, induced in a conductor, is always in such
a direction that the current produced would oppose the change that caused it; this law is a form of the law of conservation of
energy.
resonance: The increase in the amplitude of an oscillation of a system under the influence of a periodic force whose frequency
is close to that of the system’s natural frequency.
rms: Root mean square: a statistical measure of the magnitude of a varying quantity.
phasor: A representation of a complex number in terms of a complex exponential.
reactance: The opposition to the change in flow of current in an alternating current circuit, due to inductance and capacitance;
the imaginary part of the impedance.
impedance: A measure of the opposition to the flow of an alternating current in a circuit; the aggregation of its resistance,
inductive and capacitive reactance. Represented by the symbol Z.
Ohm’s law: Ohm’s observation is that the direct current flowing in an electrical circuit consisting only of resistances is directly
proportional to the voltage applied.
22.2.14 https://phys.libretexts.org/@go/page/15656
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22.3: Applications of Induction and EM Waves
learning objectives
Discuss the use of inductance in modern devices
Inductors have many uses in modern electronics. In a sound system, sound can be transmitted from a microphone to a speaker
(shown in ). The microphone works by induction, as the vibrating membrane induces an emf in a coil. That “signal” is then
transmitted to an amplifier and then to a speaker. The speaker is then driven by modulated electrical currents (produced by an
amplifier) that pass through and magnetize (by inductance) a speaker coil of copper wire, creating a magnetic field. Thus, the
electrical current variations that pass through the speaker are converted to varying magnetic forces, which move the speaker
diaphragm, forcing the driver to produce air motion that is similar to the original signal from the amplifier.
image
A simple, modern speaker.: A speaker with magnet and coils that is used to produce sound.
Inductance in modern electronics is also used in computer memory. Magnetic storage uses different patterns of magnetization on a
magnetically coated surface to store information. Differently magnetized areas on tape (or disk) induce signals on “read-write”
heads, from which the information is later accessed.
Another application is the seismograph—an instrument for detecting and recording the intensity, direction and duration of ground
movement. It contains a fixed coil and a magnet hung on a spring (or vice versa). These “record” the current induced when the
Earth shakes.
A ground fault circuit interrupter (GFCI) provides additional safety (that circuit breakers cannot) by stopping the current in a
shorted-out circuit. This is done by inductance. If a GFCI detects that there is a leakage of current, it produces an EMF and a
current in the opposite direction of the original current.
Antennae
An antenna is a device that converts electric power into radio waves, and vice versa.
learning objectives
Describe functions and uses of antennae
Maxwell’s equations predicted that all light waves have the same structure, regardless of wavelength and frequency. As a
consequence, visible light and radio waves should share common characteristics. Maxwell’s 1865 prediction passed an important
test in 1888, when Heinrich Hertz published the results of experiments in which he showed that radio waves could be manipulated
in the same ways as visible light waves. To aid in his experiment, Hertz built the first antenna.
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Car Antenna: A common car antenna that converts electric power in the air into electromagnetic waves.
An antenna (or aerial) is an electrical device that converts electric power into radio waves, and vice versa. Usually, it is used with a
radio transmitter or radio receiver. In transmission, a radio transmitter supplies an oscillating radio frequency electric current to the
antenna’s terminals, and the antenna radiates the energy from the current as electromagnetic waves (radio waves). In reception, an
antenna intercepts some of the power of an electromagnetic wave in order to produce a tiny voltage at its terminals. This voltage is
applied to a receiver to be amplified.
Antennas are essential components of all types of equipment that utilize radio. These include: radio broadcasting, broadcast
television, two-way radio, communications receivers, radar, cell phones, and satellite communications; as well as other devices
such as garage door openers, wireless microphones, bluetooth enabled devices, wireless computer networks, baby monitors, and
RFID tags on merchandise.
Antennas may also include reflective or directive elements or surfaces not connected to the transmitter or receiver, such as parasitic
elements, parabolic reflectors, or horns. These serve to direct the radio waves into a beam or other desired radiation pattern.
Antennas can be designed to transmit or receive radio waves in all directions equally (omnidirectional antennas), or transmit them
in a beam in a particular direction and receive from that one direction only (directional or high gain antennas).
Key Points
The microphone works by induction, as the vibrating membrane induces an emf in a coil.
A speaker produces sound by induction as varying magnetic forces move a speaker diaphragm, producing air motion that
produces sound.
Computer memory is stored by inducing signals on read/write heads.
GFCIs induce a current to oppose sudden increases in current in a circuit, thus preventing possible electrocution.
The first antennas were built in 1888 by German physicist Heinrich Hertz.
Antennas are essential components of all equipment that uses radio.
Antennas transmit or receive radio waves in all directions equally, or transmit them in a beam in a particular direction.
Key Terms
amplifier: An appliance or circuit that increases the strength of a weak electrical signal without changing the other
characteristics of the signal.
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ, and corresponds to the velocity of the wave divided by its frequency.
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time t in which it occurs: f = n / t.
transmitter: An electronic device that generates and amplifies a carrier wave, modulates it with a meaningful signal derived
from speech, music, TV or other sources, and broadcasts the resulting signal from an antenna.
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22.4: Magnetic Fields and Maxwell Revisited
learning objectives
Describe behavior of an inductor when the current is changed, and express energy stored in a magnetic field in a form of
equation
When a conductor carries a current, a magnetic field surrounding the conductor is produced. The resulting magnetic flux is
proportional to the current. If the current changes, the change in magnetic flux is proportional to the time-rate of change in current
by a factor called inductance (L). Since nature abhors rapid change, a voltage (electromotive force , EMF) produced in the
conductor opposes the change in current, which is also proportional to the change in magnetic flux. Thus, inductors oppose change
in current by producing a voltage that,in turn, creates a current to oppose the change in magnetic flux; the voltage is proportional to
the change in current.
(Eq. 1), where L is the inductance in units of Henry and I is the current in units of Ampere.
Magnetic Field Created By A Solenoid: Magnetic field created by a solenoid (cross-sectional view) described using field lines.
Energy is “stored” in the magnetic field.
2 2 2 2
1 1 μN A B L B
2
E = LI = = AL (22.4.2)
2
2 2 L μ2 N 2μ
2
μI A
(We used the relation L = 1
and B=μNI/LB=μNI/L. )
Therefore, the energy density uB = energy/volume of a magnetic field B is written as B = μNI/Lμ.
learning objectives
Explain how Maxwell’s prediction of the electromagnetic force was confirmed by Hertz
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1. Electric field lines originate on positive charges and terminate on negative charges, and the electric field is defined as the force
per unit charge on a test charge. The strength of the force is related to the electric constant ε0, also known as the permitivity of
free space.
2. Magnetic field lines are continuous, having no beginning or end. No magnetic monopoles are known to exist.
3. A changing magnetic field induces an electromotive force (emf) and, hence, an electric field. The direction of the emf opposes
the change.
4. Magnetic fields are generated by moving charges or by changing electric fields.
Maxwell’s equations predict that regardless of wavelength and frequency, every light wave has the same structure. This means
Maxwell’s equations predicted that radio and x-ray waves existed, even though they hadn’t actually been discovered yet.
Apparatus Used by Hertz: The apparatus used by Hertz in 1887 to generate and detect electromagnetic waves. An RLC circuit
connected to the first loop caused sparks across a gap in the wire loop and generated electromagnetic waves. Sparks across a gap in
the second loop located across the laboratory gave evidence that the waves had been received.
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EM Wave: The propogation of an electromagnetic wave as predicted by Maxwell and confirmed by Hertz.
Key Points
The formula for the energy stored in a magnetic field is E = 1/2 LI2.
The energy stored in a magnetic field is equal to the work needed to produce a current through the inductor.
2
2μ
Key Terms
inductor: A device or circuit component that exhibits significant self-inductance; a device which stores energy in a magnetic
field.
electric field: A region of space around a charged particle, or between two voltages; it exerts a force on charged objects in its
vicinity.
magnetic field: A condition in the space around a magnet or electric current in which there is a detectable magnetic force, and
where two magnetic poles are present.
electromotive force: (EMF)—The voltage generated by a battery or by the magnetic force according to Faraday’s Law. It is
measured in units of volts, not newtons, and thus, is not actually a force.
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CHAPTER OVERVIEW
23: Electromagnetic Waves is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
23.1: The Electromagnetic Spectrum
learning objectives
Compare properties of AM and FM radio waves
Radio Waves
Radio waves are a type of electromagnetic (EM) radiation with wavelengths in the electromagnetic spectrum longer than infrared
light. They have have frequencies from 300 GHz to as low as 3 kHz, and corresponding wavelengths from 1 millimeter to 100
kilometers. Like all other electromagnetic waves, radio waves travel at the speed of light. Naturally occurring radio waves are made
by lightning or by astronomical objects. Artificially generated radio waves are used for fixed and mobile radio communication,
broadcasting, radar and other navigation systems, communications satellites, computer networks and innumerable other
applications. Different frequencies of radio waves have different propagation characteristics in the Earth’s atmosphere—long waves
may cover a part of the Earth very consistently, shorter waves can reflect off the ionosphere and travel around the world, and much
shorter wavelengths bend or reflect very little and travel on a line of sight.
Electromagnetic Spectrum: The electromagnetic spectrum, showing the major categories of electromagnetic waves. The range of
frequencies and wavelengths is remarkable. The dividing line between some categories is distinct, whereas other categories
overlap. Microwaves encompass the high frequency portion of the radio section of the EM spectrum.
AM Radio Waves
AM radio waves are used to carry commercial radio signals in the frequency range from 540 to 1600 kHz. The abbreviation AM
stands for amplitude modulation—the method for placing information on these waves. A carrier wave having the basic frequency
of the radio station (for instance, 1530 kHz) is varied or modulated in amplitude by an audio signal. The resulting wave has a
constant frequency, but a varying amplitude.
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AM Radio: Amplitude modulation for AM radio. (a) A carrier wave at the station’s basic frequency. (b) An audio signal at much
lower audible frequencies. (c) The amplitude of the carrier is modulated by the audio signal without changing its basic frequency.
FM Radio Waves
FM radio waves are also used for commercial radio transmission, but in the frequency range of 88 to 108 MHz. FM stands for
frequency modulation, another method of carrying information. In this case, a carrier wave having the basic frequency of the radio
station (perhaps 105.1 MHz) is modulated in frequency by the audio signal, producing a wave of constant amplitude but varying
frequency.
FM Radio: Frequency modulation for FM radio. (a) A carrier wave at the station’s basic frequency. (b) An audio signal at much
lower audible frequencies. (c) The frequency of the carrier is modulated by the audio signal without changing its amplitude.
Since audible frequencies range up to 20 kHz (or 0.020 MHz) at most, the frequency of the FM radio wave can vary from the
carrier by as much as 0.020 MHz. For this reason, the carrier frequencies of two different radio stations cannot be closer than 0.020
MHz. An FM receiver is tuned to resonate at the carrier frequency and has circuitry that responds to variations in frequency,
reproducing the audio information.
FM radio is inherently less subject to noise from stray radio sources than AM radio because amplitudes of waves add noise. Thus,
an AM receiver would interpret noise added onto the amplitude of its carrier wave as part of the information. An FM receiver can
be fashioned to reject amplitudes other than that of the basic carrier wave and only look for variations in frequency. Thus, since
noise produces a variation in amplitude, it is easier to reject noise from FM.
TV
Electromagnetic waves also broadcast television transmission. However, as the waves must carry a great deal of visual as well as
audio information, each channel requires a larger range of frequencies than simple radio transmission. TV channels utilize
frequencies in the range of 54 to 88 MHz and 174 to 222 MHz (the entire FM radio band lies between channels 88 MHz and 174
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MHz). These TV channels are called VHF (very high frequency). Other channels called UHF (ultra high frequency) utilize an even
higher frequency range of 470 to 1000 MHz.
The TV video signal is AM, while the TV audio is FM. Note that these frequencies are those of free transmission with the user
utilizing an old-fashioned roof antenna. Satellite dishes and cable transmission of TV occurs at significantly higher frequencies,
and is rapidly evolving with the use of the high-definition or HD format.
Microwaves
Microwaves are electromagnetic waves with wavelengths ranging from one meter to one millimeter (frequencies between 300 MHz
and 300 GHz).
Learning objectives
Distinguish three ranges of the microwave portion of the electromagnetic spectrum
Microwaves
Microwaves are electromagnetic waves with wavelengths ranging from as long as one meter to as short as one millimeter, or
equivalently with frequencies between 300 MHz (0.3 GHz) and 300 GHz. The microwave region of the electromagnetic (EM)
spectrum is generally considered to overlap with the highest frequency (shortest wavelength) radio waves. As is the case for all EM
waves, microwaves travel in a vacuum at the speed of light. The prefix “micro-” in “microwave” is not meant to suggest a
wavelength in the micrometer range. It indicates that microwaves are “small” because have shorter wavelengths as compared to
waves used in typical radio broadcasting. The boundaries between far infrared light, terahertz radiation, microwaves, and ultra-
high-frequency radio waves are fairly arbitrary. They are used variously between different fields of study (see figure).
Electromagnetic Spectrum: The electromagnetic spectrum, showing the major categories of electromagnetic waves. The range of
frequencies and wavelengths is remarkable. The dividing line between some categories is distinct, whereas other categories
overlap. Microwaves overlap with the high frequency portion of the radio section of the EM spectrum.
Subcategories of Microwaves
The microwave portion of the radio spectrum can be subdivided into three ranges, listed below from high to low frequencies.
Extremely high frequency (EHF) is the highest microwave frequency band. EHF runs the range of frequencies from 30 to 300
gigahertz, above which electromagnetic radiation is considered as far infrared light, also referred to as terahertz radiation. This
frequency range corresponds to a wavelength range of 10 to 1 millimeter, so it is sometimes called the millimeter band. This
band is commonly used in radio astronomy and remote sensing.
Super high frequency (SHF) is the designation for electromagnetic wave frequencies in the range of 3 GHz to 30 GHz. This
band of frequencies is known also as the centimeter band because the wavelengths range from ten to one centimeters. This
frequency range is used for most radar transmitters, microwave ovens, wireless LANs, cell phones, satellite communication,
microwave radio relay links, and numerous short range terrestrial data links.
Ultra-high frequency (UHF) designates the microwave frequency range of electromagnetic waves between 300 MHz and 3
GHz, also known as the decimeter band because the wavelengths range from one to ten decimeters, or 10 centimeters to 1
meter. They are used for television broadcasting, cordless phones, walkie-talkies, satellite communication, and numerous other
applications.
Sources of Microwaves
Microwaves are the highest-frequency electromagnetic waves that can be produced by currents in macroscopic circuits and devices.
Microwaves can also be produced by atoms and molecules—e.g., they are a component of electromagnetic radiation generated by
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thermal agitation. The thermal motion of atoms and molecules in any object at a temperature above absolute zero causes them to
emit and absorb radiation.
Since it is possible to carry more information per unit time on high frequencies, microwaves are quite suitable for communications
devices. Most satellite-transmitted information is carried on microwaves, as are land-based long-distance transmissions. A clear
line of sight between transmitter and receiver is needed because of the short wavelengths involved.
The sun also emits microwave radiation, although most of it is blocked by Earth’s atmosphere. The Cosmic Microwave
Background Radiation (CMBR) is microwave radiation that permeates all of space, and its discovery supports the Big Bang theory
of the origin of the universe.
Cosmic Microwave Background: Cosmic background radiation of the Big Bang mapped with increasing resolution.
Cavity Magnetron: Cutaway view inside a cavity magnetron as used in a microwave oven.
Microwaves are used by microwave ovens to heat food. Microwaves at a frequency of 2.45 GHz are produced by accelerating
electrons. The microwaves then induce an alternating electric field in the oven. Water and some other constituents of food have a
slightly negative charge at one end and a slightly positive charge at one end (called polar molecules). The range of microwave
frequencies is specially selected so that the polar molecules, in trying to maintain their orientation with the electric field, absorb
these energies and increase their temperatures—a process called dielectric heating.
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Radar, first developed in World War II, is a common application of microwaves. By detecting and timing microwave echoes, radar
systems can determine the distance to objects as diverse as clouds and aircraft. A Doppler shift in the radar echo can determine the
speed of a car or the intensity of a rainstorm. Sophisticated radar systems can map the Earth and other planets, with a resolution
limited by wavelength. The shorter the wavelength of any probe, the smaller the detail it is possible to observe.
A maser is a device similar to a laser, which amplifies light energy by stimulating photons. The maser, rather than amplifying
visible light energy, amplifies the lower-frequency, longer-wavelength microwaves and radio frequency emissions.
Infrared Waves
Infrared (IR) light is EM radiation with wavelengths longer than those of visible light from 0.74 µm to 1 mm (300 GHz to 1 THz).
Skills to Develop
Distinguish three ranges of the infrared portion of the spectrum, and describe processes of absorption and emission of
infrared light by molecules
Infrared Waves
Infrared (IR) light is electromagnetic radiation with longer wavelengths than those of visible light, extending from the nominal red
edge of the visible spectrum at 0.74 micrometers (µm) to 1 mm. This range of wavelengths corresponds to a frequency range of
approximately 300 GHz to 400 THz, and includes most of the thermal radiation emitted by objects near room temperature. Infrared
light is emitted or absorbed by molecules when they change their rotational-vibrational movements.
Electromagnetic Spectrum: The electromagnetic spectrum, showing the major categories of electromagnetic waves. The range of
frequencies and wavelengths is remarkable. The dividing line between some categories is distinct, whereas other categories
overlap. Microwaves encompass the high frequency portion of the radio section of the EM spectrum.
Subcategories of IR Waves
The infrared part of the electromagnetic spectrum covers the range from roughly 300 GHz (1 mm) to 400 THz (750 nm). It can be
divided into three parts: It can be divided into three parts:
Atmospheric Transmittance: This is a plot of Earth’s atmospheric transmittance (or opacity) to various wavelengths of
electromagnetic radiation. Most UV wavelengths are absorbed by oxygen and ozone in Earth’s atmosphere. Observations of
astronomical UV sources must be done from space.
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Far-infrared, from 300 GHz (1 mm) to 30 THz (10 μm) – The lower part of this range may also be called microwaves. This
radiation is typically absorbed by so-called rotational modes in gas-phase molecules, by molecular motions in liquids, and by
phonons in solids. The water in Earth’s atmosphere absorbs so strongly in this range that it renders the atmosphere in effect
opaque. However, there are certain wavelength ranges (“windows”) within the opaque range that allow partial transmission, and
can be used for astronomy. The wavelength range from approximately 200 μm up to a few mm is often referred to as “sub-
millimeter” in astronomy, reserving far infrared for wavelengths below 200 μm.
Mid-infrared, from 30 to 120 THz (10 to 2.5 μm) – Hot objects (black-body radiators) can radiate strongly in this range, and
human skin at normal body temperature radiates strongly at the lower end of this region. This radiation is absorbed by
molecular vibrations, where the different atoms in a molecule vibrate around their equilibrium positions. This range is
sometimes called the fingerprint region, since the mid-infrared absorption spectrum of a compound is very specific for that
compound.
Near-infrared, from 120 to 400 THz (2,500 to 750 nm) – Physical processes that are relevant for this range are similar to those
for visible light. The highest frequences in this region can be detected directly by some types of photographic film, and by many
types of solid state image sensors for infrared photography and videography.
Note that in some fields the boundaries of these categories differ slightly; for example, in astronomy “near-infrared” is considered
to extend to 5 μm rather than 2.5 μm.
Sources of IR Waves
As stated above, while infrared radiation is commonly referred to as heat radiation, only objects emitting with a certain range of
temperatures and emissivities will produce most of their electromagnetic emission in the infrared part of the spectrum. However,
this is the case for most objects and environments humans encounter in our daily lives. Humans, their surroundings, and the Earth
itself emit most of their thermal radiation at wavelengths near 10 microns, the boundary between mid and far infrared according to
the delineation above. The range of wavelengths most relevant to thermally emitting objects on earth is often called the thermal
infrared. Many astronomical objects emit detectable amounts of IR radiation at non-thermal wavelengths.
Infrared radiation can be used to remotely determine the temperature of objects (if the emissivity is known). This is termed
thermography, mainly used in military and industrial applications but the technology is reaching the public market in the form of
infrared cameras on cars due to the massively reduced production costs.
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Applications of IR waves extend to heating, communication, meteorology, spectroscopy, astronomy, biological and medical
science, and even the analysis of works of art.
Visible Light
Visible light is the portion of the electromagnetic spectrum that is visible to the human eye, ranging from roughly 390 to 750 nm.
Skills to Develop
Distinguish six ranges of the visible spectrum
Visible Light
Visible light, as called the visible spectrum, is the portion of the electromagnetic spectrum that is visible to (can be detected by) the
human eye. Electromagnetic radiation in this range of wavelengths is often simply referred to as “light”. A typical human eye will
respond to wavelengths from about 390 to 750 nm (0.39 to 0.75 µm). In terms of frequency, this corresponds to a band in the
vicinity of 400–790 THz. A light-adapted eye generally has its maximum sensitivity at around 555 nm (540 THz), in the green
region of the optical spectrum. The spectrum does not, however, contain all the colors that the human eyes and brain can
distinguish. Unsaturated colors such as pink, or purple variations such as magenta, are absent, for example, because they can be
made only by a mix of multiple wavelengths.
Electromagnetic Spectrum: The electromagnetic spectrum, showing the major categories of electromagnetic waves. The range of
frequencies and wavelengths is remarkable. The dividing line between some categories is distinct, whereas other categories
overlap. Microwaves encompass the high frequency portion of the radio section of the EM spectrum.
Visible light is produced by vibrations and rotations of atoms and molecules, as well as by electronic transitions within atoms and
molecules. The receivers or detectors of light largely utilize electronic transitions. We say the atoms and molecules are excited
when they absorb and relax when they emit through electronic transitions.
Visible Spectrum: A small part of the electromagnetic spectrum that includes its visible components. The divisions between
infrared, visible, and ultraviolet are not perfectly distinct, nor are those between the seven rainbow colors.
The figure above shows this part of the spectrum, together with the colors associated with particular pure wavelengths. Red light
has the lowest frequencies and longest wavelengths, while violet has the highest frequencies and shortest wavelengths. Blackbody
radiation from the Sun peaks in the visible part of the spectrum but is more intense in the red than in the violet, making the Sun
yellowish in appearance.
Colors that can be produced by visible light of a narrow band of wavelengths (monochromaticlight) are called pure spectral colors.
Quantitatively, the regions of the visible spectrum encompassing each spectral color can be delineated roughly as:
red – 620 to 750 nm (400-484 THz)
Note that each color can come in many shades, since the spectrum is continuous. The human eye is insensitive to electromagnetic
radiation outside this range. By definition any images presented with data recorded from wavelengths other than those in the visible
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part of the spectrum (such as IR images of humans or animals or astronomical X-ray images) are necessarily in false color.
Atmospheric Transmittance: This is a plot of Earth’s atmospheric transmittance (or opacity) to various wavelengths of
electromagnetic radiation. Most UV wavelengths are absorbed by oxygen and ozone in Earth’s atmosphere. Observations of
astronomical UV sources must be done from space.
The optical window is also called the visible window because it overlaps the human visible response spectrum. This is not
coincidental as humanity’s ancestors evolved vision that could make use of the most plentiful wavelengths of light.The near
infrared (NIR) window lies just out of the human vision, as well as the Medium Wavelength IR (MWIR) window and the Long
Wavelength or Far Infrared (LWIR or FIR) window though other animals may experience them.
A consequence of the existence of the optical window in Earth’s atmosphere is the relatively balmy temperature conditions on
Earth’s surface. The Sun’s luminosity function peaks in the visible range and light in that range is able to travel to the surface of the
planet unattenuated due to the optical window. This allows visible light to heat the surface. The surface of the planet then emits
energy primarily in infrared wavelengths, which has much greater difficulty escaping (and thus causing the planet to cool) due to
the opacity of the atmosphere in the infrared. Earth’s surface would be much cooler without this effect.
Photosynthesis
Plants, like animals, have evolved to utilize and respond to parts of the electromagnetic spectrum they are embedded in. Plants (and
many bacteria) convert the light energy captured from the Sun into chemical energy that can be used to fuel the organism’s
activities. In plants, algae, and cyanobacteria, photosynthesis uses carbon dioxide and water, releasing oxygen as a waste product.
Photosynthesis is vital for all aerobic life on Earth (such as humans and animals). The portion of the EM spectrum used by
photosynthesic organisms is called the photosynthetically active region (PAR) and corresponds to solar radiation between 400 and
700 nm, substantially overlapping with the range of human vision. This is again not coincidental; the light in this range is the most
plentiful to organisms on the surface of Earth because the Sun emits about half of its luminosity in this wavelength range and it is
allowed to pass freely through the optical windows in Earth’s atmosphere.
Ultraviolet Light
Ultraviolet (UV) light is electromagnetic radiation with a wavelength shorter than that of visible light in the range 10 nm to 400
nm.
Skills to Develop
Identify wavelength range characteristic for ultraviolet light and its biological effects
Ultraviolet Light
Ultraviolet (UV) light is electromagnetic radiation with a wavelength shorter than that of visible light, but longer than X-rays, that
is, in the range 10 nm to 400 nm, corresponding to photon energies from 3 eV to 124 eV (1 eV = 1.6e-19 J; EM radiation with
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frequencies higher than those of visible light are often expressed in terms of energy rather than frequency). It is so-named because
the spectrum consists of electromagnetic waves with frequencies higher than those that humans identify as the color violet. These
frequencies are invisible to humans, but visible to a number of insects and birds.
Electromagnetic Spectrum: The electromagnetic spectrum, showing the major categories of electromagnetic waves. The range of
frequencies and wavelengths is remarkable. The dividing line between some categories is distinct, whereas other categories
overlap. Microwaves encompass the high frequency portion of the radio section of the EM spectrum.
UV light is found in sunlight (where it constitutes about 10% of the energy in vacuum) and is emitted by electric arcs and
specialized lights such as black lights. It can cause chemical reactions, and causes many substances to glow or fluoresce. Most
ultraviolet is classified as non-ionizing radiation. The higher energies of the ultraviolet spectrum from wavelengths about 10 nm to
120 nm (‘extreme’ ultraviolet) are ionizing, but this type of ultraviolet in sunlight is blocked by normal molecular oxygen (O2) in
air, and does not reach the ground. However, the entire spectrum of ultraviolet radiation has some of the biological features of
ionizing radiation, in doing far more damage to many molecules in biological systems than is accounted for by simple heating
effects (an example is sunburn). These properties derive from the ultraviolet photon’s power to alter chemical bonds in molecules,
even without having enough energy to ionize atoms.
Although ultraviolet radiation is invisible to the human eye, most people are aware of the effects of UV on the skin, called suntan
and sunburn. In addition to short wave UV blocked by oxygen, a great deal (>97%) of mid-range ultraviolet (almost all UV above
280 nm and most up to 315 nm) is blocked by the ozone layer, and like ionizing short wave UV, would cause much damage to
living organisms if it penetrated the atmosphere. After atmospheric filtering, only about 3% of the total energy of sunlight at the
zenith is ultraviolet, and this fraction decreases at other sun angles. Much of it is near-ultraviolet that does not cause sunburn, but is
still capable of causing long term skin damage and cancer. An even smaller fraction of ultraviolet that reaches the ground is
responsible for sunburn and also the formation of vitamin D (peak production occurring between 295 and 297 nm) in all organisms
that make this vitamin (including humans). The UV spectrum thus has many effects, both beneficial and damaging, to human
health.
Atmospheric Transmittance: This is a plot of Earth’s atmospheric opacity (opposite of transmittance) to various wavelengths of
electromagnetic radiation, including visible light. Visible light passes relatively unimpeded through the atmosphere in the “optical
window.” Most UV wavelengths are absorbed by oxygen and ozone in Earth’s atmosphere. Observations of astronomical UV
sources must be done from space.
Subcategories of UV Light
Solar UV radiation is commonly subdivided into three regions: UV-A (320–400 nm), UV-B (290–320 nm), and UV-C (220–290
nm), ranked from long to shorter wavelengths (from smaller to larger energies). Most UV-B and all UV-C is absorbed by ozone
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(O3) molecules in the upper atmosphere. Consequently, 99% of the solar UV radiation reaching the Earth’s surface is UV-A.
There are other schemes for dividing UV into different categories, another common one is: near-ultraviolet (NUV – 300-400 nm),
middle ultraviolet (MUV – 200-300 nm), far ultraviolet (FUV – 200-122 nm), and extreme ultraviolet (EUV- 121-10 nm).
Harmful Effects
An overexposure to UVB radiation can cause sunburn and some forms of skin cancer. In humans, prolonged exposure to solar UV
radiation may result in acute and chronic health effects on the skin, eye, and immune system. Moreover, UVC can cause adverse
effects that can variously be mutagenic or carcinogenic.
DNA UV Mutation: Ultraviolet photons harm the DNA molecules of living organisms in different ways. In one common damage
event, adjacent thymine bases bond with each other, instead of across the “ladder. ” This “thymine dimer” makes a bulge, and the
distorted DNA molecule does not function properly.
The International Agency for Research on Cancer of the World Health Organization has classified all categories and wavelengths
of ultraviolet radiation as a Group 1 carcinogen. This is the highest level designation for carcinogens and means that “there is
enough evidence to conclude that it can cause cancer in humans. ”
Beneficial Effects
UVB exposure induces the production of vitamin D in the skin. The majority of positive health effects are related to this vitamin. It
has regulatory roles in calcium metabolism (which is vital for normal functioning of the nervous system, as well as for bone growth
and maintenance of bone density), immunity, cell proliferation, insulin secretion, and blood pressure.
X-Rays
X-rays are electromagnetic waves with wavelengths in the range of 0.01 to 10 nanometers and energies in the range of 100 eV to
100 keV.
Skills to Develop
Distinguish two categories of X-rays and their biological effects
X-Rays
X-rays are electromagnetic waves with wavelengths in the range of 0.01 to 10 nanometers, corresponding to frequencies in the
range 30 petahertz to 30 exahertz (3×1016 Hz to 3×1019 Hz) and energies in the range 100 eV to 100 keV. They are shorter in
wavelength than UV rays and longer than gamma rays. In many languages, X-radiation is called Röntgen radiation, after Wilhelm
Röntgen, who is usually credited as its discoverer, and who had named it X-radiation to signify an unknown type of radiation.
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Electromagnetic Spectrum: The electromagnetic spectrum, showing the major categories of electromagnetic waves. The range of
frequencies and wavelengths is remarkable. The dividing line between some categories is distinct, whereas other categories
overlap. Microwaves encompass the high frequency portion of the radio section of the EM spectrum.
X-Ray Spectrum and Applications: X-rays are part of the electromagnetic spectrum, with wavelengths shorter than those of
visible light. Different applications use different parts of the X-ray spectrum.
X-rays with photon energies above 5 to 10 keV (below 0.2-0.1 nm wavelength), are called hard X-rays, while those with lower
energy are called soft X-rays. Due to their penetrating ability, hard X-rays are widely used to image the inside of objects (e.g., in
medical radiography and airport security). As a result, the term X-ray is metonymically used to refer to a radiographic image
produced using this method, in addition to the method itself. Since the wavelength of hard X-rays are similar to the size of atoms,
they are also useful for determining crystal structures by X-ray crystallography. By contrast, soft X-rays are easily absorbed in air
and the attenuation length of 600 eV (~2 nm) X-rays in water is less than 1 micrometer.
In medical diagnostic applications, the low energy (soft) X-rays are unwanted, since they are totally absorbed by the body,
increasing the radiation dose without contributing to the image. Hence, a thin metal sheet, often of aluminum, called an X-ray filter,
is usually placed over the window of the X-ray tube, absorbing the low energy part in the spectrum. This is called hardening the
beam since it shifts the center of the spectrum towards higher energy (or harder) X-rays.
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Gamma Rays
Gamma rays are very high frequency electromagnetic waves usually emitted from radioactive decay with frequencies greater than
1019Hz.
Skills to Develop
Identify wavelength range characteristic for gamma rays, noting their biological effects and distinguishing them from
gamma rays
Gamma Rays
Gamma radiation, also known as gamma rays or hyphenated as gamma-rays and denoted as γ, is electromagnetic radiation of high
frequency and therefore high energy. Gamma rays typically have frequencies above 10 exahertz (or >1019 Hz), and therefore have
energies above 100 keV and wavelengths less than 10 picometers (less than the diameter of an atom). However, this is not a hard
and fast definition, but rather only a rule-of-thumb description for natural processes. Gamma rays from radioactive decay are
defined as gamma rays no matter what their energy, so that there is no lower limit to gamma energy derived from radioactive decay.
Gamma decay commonly produces energies of a few hundred keV, and almost always less than 10 MeV.
Gamma Decay: Illustration of an emission of a gamma ray (γ) from an atomic nucleus
Gamma rays are ionizing radiation and are thus biologically hazardous. They are classically produced by the decay from high
energy states of atomic nuclei, a process called gamma decay, but are also created by other processes. Paul Villard, a French
chemist and physicist, discovered gamma radiation in 1900, while studying radiation emitted from radium during its gamma decay.
Villard’s radiation was named “gamma rays” by Ernest Rutherford in 1903.
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Exceptions to this convention occur in astronomy, where gamma decay is seen in the afterglow of certain supernovas, but other
high energy processes known to involve other than radioactive decay are still classed as sources of gamma radiation. A notable
example is extremely powerful bursts of high-energy radiation normally referred to as long duration gamma-ray bursts, which
produce gamma rays by a mechanism not compatible with radioactive decay. These bursts of gamma rays, thought to be due to the
collapse of stars called hypernovas, are the most powerful events so far discovered in the cosmos. Astrophysical processes are the
only sources for very high energy gamma rays (~100 MeV).
Gamma Ray Sky Map: This is an image of the entire sky in 100 MeV or greater gamma rays as seen by the EGRET instrument
aboard the CGRO spacecraft. Bright spots within the galactic plane are pulsars (spinning neutron stars with strong magnetic fields),
while those above and below the plane are thought to be quasars (galaxies with supermassive black holes actively accreting matter).
Health Effects
All ionizing radiation causes similar damage at a cellular level, but because rays of alpha particles and beta particles are relatively
non-penetrating, external exposure to them causes only localized damage (e.g., radiation burns to the skin). Gamma rays and
neutrons are more penetrating, causing diffuse damage throughout the body (e.g., radiation sickness, cell’s DNA damage, cell death
due to damaged DNA, increasing incidence of cancer) rather than burns. External radiation exposure should also be distinguished
from internal exposure, due to ingested or inhaled radioactive substances, which, depending on the substance’s chemical nature,
can produce both diffuse and localized internal damage. The most biological damaging forms of gamma radiation occur at energies
between 3 and 10 MeV.
Key Points
The lowest frequency portion of the electromagnetic spectrum is designated as “radio,” generally considered to have
wavelengths within 1 millimeter to 100 kilometers or frequencies within 300 GHz to 3 kHz.
There is a wide range of subcategories contained within radio including AM and FM radio. Radio waves can be generated by
natural sources such as lightning or astronomical phenomena; or by artificial sources such as broadcast radio towers, cell
phones, satellites and radar.
AM radio waves are used to carry commercial radio signals in the frequency range from 540 to 1600 kHz. The abbreviation
AM stands for amplitude modulation—the method for placing information on these waves. AM waves have constant frequency,
but a varying amplitude.
FM radio waves are also used for commercial radio transmission in the frequency range of 88 to 108 MHz. FM stands for
frequency modulation, which produces a wave of constant amplitude but varying frequency.
The microwave region of the electromagnetic (EM) spectrum is generally considered to overlap with the highest frequency
(shortest wavelength ) radio waves.
The prefix “micro-” in “microwave” is not meant to suggest a wavelength in the micrometer range. It indicates that microwaves
are “small” compared to waves used in typical radio broadcasting in that they have shorter wavelengths.
The microwave portion of the electromagnetic spectrum can be subdivided into three ranges listed below from high to low
frequencies: extremely high frequency (30 to 300 GHz), super high frequency (3 to 30 GHz), and ultra-high frequency (300
MHz to 3 GHz).
Microwave sources include artificial devices such as circuits, transmission towers, radar, masers, and microwave ovens, as well
as natural sources such as the Sun and the Cosmic Microwave Background.
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Microwaves can also be produced by atoms and molecules. They are, for example, a component of electromagnetic radiation
generated by thermal agitation. The thermal motion of atoms and molecules in any object at a temperature above absolute zero
causes them to emit and absorb radiation.
Infrared light includes most of the thermal radiation emitted by objects near room temperature. Infrared light is emitted or
absorbed by molecules when they change their rotational-vibrational movements.
The infrared portion of the spectrum can be divided into three regions in wavelength: far-infrared, from 300 GHz (1 mm) to 30
THz (10 μm); mid-infrared, from 30 to 120 THz (10 to 2.5 μm); and near-infrared, from 120 to 400 THz (2,500 to 750 nm).
Infrared radiation is popularly known as ” heat radiation,” but light and electromagnetic waves of any frequency will heat
surfaces that absorb them.
The concept of emissivity is important in understanding the infrared emissions of objects. This is a property of a surface which
describes how its thermal emissions deviate from the ideal of a black body.
Infrared radiation can be used to remotely determine the temperature of objects (if the emissivity is known). This is termed
thermography, mainly used in military and industrial applications.
Visible light is produced by vibrations and rotations of atoms and molecules, as well as by electronic transitions within atoms
and molecules. We say the atoms and molecules are excited when they absorb and relax when they emit through electronic
transitions.
This figure shows the visible part of the spectrum, together with the colors associated with particular pure wavelengths. Red
light has the lowest frequencies and longest wavelengths, while violet has the highest frequencies and shortest wavelengths.
Colors that can be produced by visible light of a narrow band of wavelengths are called pure spectral colors. They can be
delineated roughly in wavelength as: violet (380-450 nm), blue (450-495 nm), green (495-570 nm), yellow (570-590 nm),
orange (590-620 nm), and red (620 to 750 nm).
Visible wavelengths pass through the optical window, the Earth’s atmosphere allows this region of the electromagnetic
spectrum to pass through largely unattenuated (see opacity plot in.
The portion of the EM spectrum used by photosynthesic organisms is called the photosynthetically active region (PAR) and
corresponds to solar radiation between 400 and 700 nm, substantially overlapping with the range of human vision.
Ultraviolet light gets its name because the spectrum consists of electromagnetic waves with frequencies higher than those that
humans identify as the color violet.
Most UV is non- ionizing radiation, though UV with higher energies (10-120 nm) is ionizing. All UV can have harmful effects
on biological matter (such as causing cancers) with the highest energies causing the most damage.
The danger posed by lower energy UV radiation is derived from the ultraviolet photon ‘s power to alter chemical bonds in
molecules, even without having enough energy to ionize atoms.
Solar UV radiation is commonly subdivided into three regions: UV-A (320–400 nm), UV-B (290–320 nm), and UV-C (220–290
nm), ranked from long to shorter wavelengths (from smaller to larger energies).
Most UV-B and all UV-C is absorbed by ozone (O3) molecules in the upper atmosphere. Consequently, 99% of the solar UV
radiation reaching the Earth’s surface is UV-A.
X-rays have shorter wavelengths (higher energy ) than UV waves and, generally, longer wavelengths (lower energy) than
gamma rays. Sometimes X-rays are called Röntgen radiation, after Wilhelm Röntgen, who is usually credited as their
discoverer.
Because X-rays have very high energy they are known as ionizing radiation and can harm living tissue. A very high radiation
dose over a short amount of time causes radiation sickness, while lower doses can give an increased risk of radiation-induced
cancer.
Lower doses of X-ray radiation can be very effectively used in medical radiography and X-ray spectroscopy. In the case of
medical radiography, the benefits of using X-rays for examination far outweighs the risk.
X-rays are broken up into broad two categories: hard X-rays with energies above 5-10 keV (below 0.2-0.1 nm wavelength) and
soft X-rays with energies 100 eV – 5 keV (10 – 0.1 nm wavelength). Hard X-rays are more useful for radiography because they
pass through tissue.
The distinction between X-rays and gamma rays is somewhat arbitrary and there is substantial overlap at the high energy
boundary. However, in general they are distinguished by their source, with gamma rays originating from the nucleus and X-rays
from the electrons in the atom.
Gamma rays are the highest energy EM radiation and typically have energies greater than 100 keV, frequencies greater than
1019 Hz, and wavelengths less than 10 picometers.
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Gamma rays from radioactive decay are defined as gamma rays no matter what their energy, so that there is no lower limit to
gamma energy derived from radioactive decay. Gamma decay commonly produces energies of a few hundred keV, and almost
always less than 10 MeV.
Gamma rays have characteristics identical to X-rays of the same frequency—they differ only in source. Gamma rays are usually
distinguished by their origin: X-rays are emitted by definition by electrons outside the nucleus, while gamma rays are emitted
by the nucleus.
Natural sources of gamma rays include gamma decay from naturally occurring radioisotopes such as potassium-40, and also as
a secondary radiation from atmospheric interactions with cosmic ray particles. Exotic astrophysical processes will also produce
gamma rays.
Gamma rays are ionizing radiation and are thus biologically hazardous. The most biological damaging forms of gamma
radiation occur at energies between 3 and 10 MeV.
Key Terms
AM radio waves: Waves used to carry commercial radio signals between 540 and 1600 kHz. Information is carried by
amplitude variation, while the frequency remains constant.
FM radio waves: Waves used to carry commercial radio signals between 88 and 108 MHz. Information is carried by frequency
modulation, while the signal amplitude remains constant.
radio waves: Designates a portion of the electromagnetic spectrum having frequencies ranging from 300 GHz to 3 kHz, or
equivalently, wavelengths from 1 millimeter to 100 kilometers.
terahertz radiation: Electromagnetic waves with frequencies around one terahertz.
thermal agitation: The thermal motion of atoms and molecules in any object at a temperature above absolute zero, causing
them to emit and absorb radiation.
radar: A method of detecting distant objects and determining their position, velocity, or other characteristics by analysis of sent
radio waves (usually microwaves) reflected from their surfaces.
emissivity: The energy-emitting propensity of a surface, usually measured at a specific wavelength.
thermography: Any of several techniques for the remote measurement of the temperature variations of a body, especially by
creating images produced by infrared radiation.
thermal radiation: The electromagnetic radiation emitted from a body as a consequence of its temperature; increasing the
temperature of the body increases the amount of radiation produced, and shifts it to shorter wavelengths (higher frequencies) in
a manner explained only by quantum mechanics.
spectral color: a color that is evoked by a single wavelength of light in the visible spectrum, or by a relatively narrow band of
wavelengths. Every wavelength of light is perceived as a spectral color, in a continuous spectrum; the colors of sufficiently
close wavelengths are indistinguishable.
optical window: the optical portion of the electromagnetic spectrum that passes through the atmosphere all the way to the
ground. The window runs from around 300 nanometers (ultraviolet-C) at the short end up into the range the eye can use,
roughly 400-700 nm and continues up through the visual infrared to around 1100 nm, which is thermal infrared.
visible light: the part of the electromagnetic spectrum, between infrared and ultraviolet, that is visible to the human eye
ozone layer: A region of the stratosphere, between 15 and 30 kilometres in altitude, containing a relatively high concentration
of ozone; it absorbs most solar ultraviolet radiation.
ionizing radiation: high-energy radiation that is capable of causing ionization in substances through which it passes; also
includes high-energy particles
non-ionizing radiation: Radiation that does not cause atmospheric ionization; electrically neutral radiation.
X-ray spectroscopy: The use of an X-ray spectrometer for chemical analysis.
x-ray crystallography: A technique in which the patterns formed by the diffraction of X-rays on passing through a crystalline
substance yield information on the lattice structure of the crystal, and the molecular structure of the substance.
radiograph: An image, often a photographic negative, produced by radiation other than normal light; especially an X-ray
photograph.
gamma ray: A very high frequency (and therefore very high energy) electromagnetic radiation emitted as a consequence of
radioactivity.
gamma decay: A nuclear reaction with the emission of a gamma ray.
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23.2: Electromagnetic Waves and their Properties
learning objectives
Explain the meaning and importance of Maxwell’s equations
Maxwell’s Equations
Maxwell’s equations are a set of four partial differential equations that, along with the Lorentz force law, form the foundation of
classical electrodynamics, classical optics, and electric circuits.
Named after esteemed physicist James Clerk Maxwell, the equations describe the creation and propagation of electric and magnetic
fields. Fundamentally, they describe how electric charges and currents create electric and magnetic fields, and how they affect each
other.
Maxwell’s equations can be divided into two major subsets. The first two, Gauss’s law and Gauss’s law for magnetism, describe
how fields emanate from charges and magnets respectively. The other two, Faraday’s law and Ampere’s law with Maxwell’s
correction, describe how induced electric and magnetic fields circulate around their respective sources.
Each of Maxwell’s equations can be looked at from the “microscopic” perspective, which deals with total charge and total current,
and the “macroscopic” set, which defines two new auxiliary fields that allow one to perform calculations without knowing
microscopic data like atomic-level charges.
Gauss’s Law
Gauss’s law relates an electric field to the charge(s) that create(s) it. The field (E) points towards negative charges and away from
positive charges, and from the microscopic perspective, is related to charge density (ρ) and vaccuum permittivity (ε0, or
permittivity of free space) as:
ρ
∇⋅E = (23.2.1)
ϵ0
Gauss’s Law basically says that a net amount of charge contained within a region of space will generate an electric field that
emanates through the surface that surrounds that region.
Example of Gauss’s Law: A positive charge contained within a region of space creates an electric field that emanates from the
surface of that region.
23.2.1 https://phys.libretexts.org/@go/page/15712
Field lines caused by a magnetic dipole: The field lines created by this magnetic dipole either form loops or extend infinitely.
The differential form of Gauss’s law for magnetic for magnetism is
∇⋅B = 0 (23.2.2)
Faraday’s Law
Faraday’s law describes how a time-varying magnetic field (or flux) induces an electric field. The principle behind this
phenomenon is used in many electric generators. Both macroscopic and microscopic differential equations are the same, relating
electric field (E) to the time-partial derivative of magnetic field (B):
∂B
∇×E = − (23.2.3)
∂t
Electromagnetic Waves: Electric (red) and magnetic (blue) waves propagate in phase sinusoidally, and perpendicularly to one
another.
23.2.2 https://phys.libretexts.org/@go/page/15712
The microscopic approach to the Maxwell-corrected Ampere’s law relates magnetic field (B) to current density (J, or current per
unit cross sectional area) and the time-partial derivative of electric field (E):
∂E
∇ × B = μ0 J + μ0 ϵ0 (23.2.4)
∂t
learning objectives
Explain the self-perpetuating behavior of an electromagnetic wave
Electromagnetic waves
Electromagnetic radiation, is a form of energy emitted by moving charged particles. As it travels through space it behaves like a
wave, and has an oscillating electric field component and an oscillating magnetic field. These waves oscillate perpendicularly to
and in phase with one another.
Electromagnetic Wave: Electromagnetic waves are a self-propagating transverse wave of oscillating electric and magnetic fields.
The direction of the electric field is indicated in blue, the magnetic field in red, and the wave propagates in the positive x-direction.
Notice that the electric and magnetic field waves are in phase.
The creation of all electromagnetic waves begins with a charged particle. This charged particle creates an electric field (which can
exert a force on other nearby charged particles). When it accelerates as part of an oscillatory motion, the charged particle creates
ripples, or oscillations, in its electric field, and also produces a magnetic field (as predicted by Maxwell’s equations).
Once in motion, the electric and magnetic fields created by a charged particle are self-perpetuating—time-dependent changes in
one field (electric or magnetic) produce the other. This means that an electric field that oscillates as a function of time will produce
a magnetic field, and a magnetic field that changes as a function of time will produce an electric field. Both electric and magnetic
fields in an electromagnetic wave will fluctuate in time, one causing the other to change.
Electromagnetic waves are ubiquitous in nature (i.e., light) and used in modern technology—AM and FM radio, cordless and
cellular phones, garage door openers, wireless networks, radar, microwave ovens, etc. These and many more such devices use
electromagnetic waves to transmit data and signals.
All the above sources of electromagnetic waves use the simple principle of moving charge, which can be easily modeled. Placing a
coin in contact with both terminals of a 9-volt battery produces electromagnetic waves that can be detected by bringing the antenna
of a radio (tuned to a static-producing station) within a few inches of the point of contact.
learning objectives
Relate energy of an electromagnetic wave with the frequency and wavelength
Electromagnetic radiation can essentially be described as photon streams. These photons are strictly defined as massless, but have
both energy and surprisingly, given their lack of mass, momentum, which can be calculated from their wave properties.
Waves were poorly understood until the 1900s, when Max Planck and Albert Einstein developed modern corrections to classical
theory.
23.2.3 https://phys.libretexts.org/@go/page/15712
Planck theorized that “black bodies” (thermal radiators) and other forms of electromagnetic radiation existed not as spectra, but in
discrete, “quantized” form. In other words, there were only certain energies an electromagnetic wave could have. In his work he
developed what is now known as “Planck’s constant,” which is approximately equal to 6.626×10-34 J·s.
Energy
The energy (E) of a photon can be related to its frequency (f) by Planck’s constant (h):
hc
E = hf = (23.2.5)
λ
The ratio of speed of light (c) to wavelength (λ) can be substituted in place of f to give the same equation to energy in different
terms. Note that energy cannot take any value: it can only exist in increments of frequency times Planck’s constant (or Planck’s
constant times c divided by wavelength). Energy of a wave is therefore “quantized. ”
Momentum
Momentum is classically defined as the product of mass and velocity and thus would intuitively seem irrelevant to a discussion of
electromagnetic radiation, which is both massless and composed of waves.
However, Einstein proved that light can act as particles in some circumstances, and that a wave-particle duality exists. And, given
that he related energy and mass (E=mc2), it becomes more conceivable that a wave (which has an energy value) not only has an
equation to mass but a momentum as well.
And indeed, Einstein proved that the momentum (p) of a photon is the ratio of its energy to the speed of light.
E hf h
p= = = (23.2.6)
c c λ
Substituting E with hc/λ cancels the c terms, making momentum also equal to the simple ratio of Planck’s constant to wavelength.
learning objectives
Relate speed of light with the index of refraction of the medium
Light Going from Earth to the Moon: A beam of light is depicted travelling between the Earth and the Moon in the time it takes
a light pulse to move between them: 1.255 seconds at their mean orbital (surface-to-surface) distance. The relative sizes and
separation of the Earth–Moon system are shown to scale.
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It is just that: the speed of a photon or light particle. The speed of light in a vacuum (commonly written as c) is 299,792,458 meters
per second. This is a universal physical constant used in many areas of physics. For example, you might be familiar with the
equation:
2
E = mc (23.2.7)
where E = Energy and m = mass. This is known as the mass-energy equivalence, and it uses the speed of light to interrelate space
and time. This not only explains the energy a body of mass contains, but also explains the hindrance mass has on speed.
There are many uses for the speed of light in a vacuum, such as in special relativity, which says that c is the natural speed limit and
nothing can move faster than it. However, we know from our understanding of physics (and previous atoms) that the speed at
which something travels also depends on the medium through which it is traveling. The speed at which light propagates through
transparent materials (air, glass, etc.,) is dependent on the refractive index of that material, n:
c
v = (23.2.8)
n
where v = actual velocity of light moving through the medium, c = speed of light in a vacuum, and n = refractive index of medium.
The refractive index of air is about 1.0003, and from this equation we can find that the speed of visible light in air is about 90 km/s
slower than c.
As mentioned earlier, the speed of light (usually of light in a vacuum) is used in many areas of physics. Below is an example of an
application of the constant c.
At low velocities, the quotient of v2/c2 is sufficiently close to 0 such that γ is approximately 1. However, as velocity approaches c, γ
increases rapidly towards infinity.
learning objectives
Give examples of daily observations of the Doppler effect
The Doppler Effect and Sirens: Waves emitted by a siren in a moving vehicle
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The difference in the perceived pitch depending on observer location can be explained by the fact that the siren’s position changes
as it emits waves. A wave of sound is emitted by a moving vehicle every millisecond. The vehicle ‘chases’ each wave in one
direction. By the time the next wave is emitted, it is closer (relative to an onlooker ahead of the vehicle) to the previous wave than
the wave’s frequency would suggest. Relative to an onlooker behind the vehicle, the second wave is further from the first wave
than one would expect, which suggests a lower frequency.
The Doppler effect can be caused by any kind of motion. In the example above, the siren moved relative to a stationary observer. If
the observer moves relative to the stationary siren, the observer will notice the Doppler effect on the pitch of the siren. Finally, if
the medium through which the waves propagate moves, the Doppler effect will be noticed even for a stationary observer. An
example of this phenomenon is wind.
Quantitatively, the Doppler effect can be characterized by relating the frequency perceived (f) to the velocity of waves in the
medium (c), the velocity of the receiver relative to the medium (vr), the velocity of the source relative to the medium (vs), and the
actual emitted frequency (f0):
c + vr
f =( ) f0 (23.2.10)
c + vs
learning objectives
Explain formation of radiation pressure
Radiation pressure is the pressure exerted upon any surface exposed to electromagnetic (EM) radiation. EM radiation (or photon,
which is a quantum of light) carries momentum; this momentum is transferred to an object when the radiation is absorbed or
reflected. Perhaps one of the most well know examples of the radiation pressure would be comet tails. Haley’s comet is shown in.
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Halley’s Comet: As a comet approaches the inner Solar System, solar radiation causes the volatile materials within the comet to
vaporize and stream out of the nucleus. The streams of dust and gas thus released form an atmosphere around the comet (called the
coma), and the force exerted on the coma by the Sun’s radiation pressure and solar wind cause the formation of an enormous tail
that points away from the Sun.
Although radiation pressure can be understood using classical electrodynamics, here we will examine the quantum mechanical
argument. From the perspective of quantum theory, light is made of photons: particles with zero mass but which carry energy and –
importantly in this argument – momentum. According to special relativity, because photons are devoid of mass, their energy (E)
and momentum (p) are related by E=pc.
Now consider a beam of light perpendicularly incident on a surface, and let us assume the beam of light is totally absorbed. The
momentum the photons carry is a conserved quantity (i.e., it cannot be destroyed) so it must be transferred to the surface; thus the
absorption of the light beam causes the surface to gain momentum. Newton’s Second Law tells us that force equals rate of change
of momentum; thus during each second, the surface experiences a force (or pressure, as pressure is force per unit area) due to the
momentum the photons transfer to it.
This gives us: pressure = momentum transferred per second per unit area = energy deposited per second per unit area / c = I/c,
(where I is the intensity of the beam of light).
Laser Cooling
There are many variations of laser cooling, but they all use radiation pressure to remove energy from atomic gases (and therefore
cool the sample). In laser cooling (sometimes called Doppler cooling), the frequency of light is tuned slightly below an electronic
transition in the atom. Because light is detuned to the “red” (i.e., at lower frequency) of the transition, the atoms will absorb more
photons if they move towards the light source, due to the Doppler effect. Thus if one applies light from two opposite directions, the
atoms will always scatter more photons from the laser beam pointing opposite to their direction of motion (typical setups applies
three opposing pairs of laser beams as in ).
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The Magneto Optical Trap: Experimental setup of Magneto Optical Trap (MOT), which uses radiation pressure to cool atomic
species. Atoms are slowed down by absorbing (and emitting) photons.
In each scattering event, the atom loses a momentum equal to the momentum of the photon. If the atom (which is now in the
excited state) then emits a photon spontaneously, it will be kicked by the same amount of momentum, only in a random direction.
Since the initial momentum loss was opposite to the direction of motion (while the subsequent momentum gain was in a random
direction), the overall result of the absorption and emission process is to reduce the speed of the atom. If the absorption and
emission are repeated many times, the average speed (and therefore the kinetic energy ) of the atom will be reduced. Since the
temperature of a group of atoms is a measure of the average random internal kinetic energy, this is equivalent to cooling the atoms.
Simple laser cooling setups can produce a cold sample of atomic gases at around 1mK (=10-3 K) starting from a room temperature
gas.
Key Points
Maxwell’s four equations describe how electric charges and currents create electric and magnetic fields, and how they affect
each other.
Gauss’s law relates an electric field to the charge(s) that create(s) it.
Gauss’s law for magnetism states that there are no “magnetic charges” analogous to electric charges, and that magnetic fields
are instead generated by magnetic dipoles.
Faraday’s law describes how a time-varying magnetic field (or flux ) induces an electric field. The principle behind this
phenomenon is used in many electric generators.
Ampere ‘s law originally stated that a magnetic field is created by an electrical current. Maxwell added that a changing electric
flux can also generate a magnetic field.
Electromagnetic waves consist of both electric and magnetic field waves. These waves oscillate in perpendicular planes with
respect to each other, and are in phase.
The creation of all electromagnetic waves begins with an oscillating charged particle, which creates oscillating electric and
magnetic fields.
Once in motion, the electric and magnetic fields that a charged particle creates are self-perpetuating: time-dependent changes in
one field (electric or magnetic) produce the other.
Max Planck proved that energy of a photon (a stream of which is an electromagnetic wave ) is quantized and can exist in
multiples of “Planck’s constant” (denoted as h, approximately equal to 6.626×10-34 J·s).
E = hf =
hc
λ
describes the energy (E) of a photon as a function of frequency (f), or wavelength (λ).
p=
E
c
=
hf
c
=
h
λ
describes the momentum (p) of a photon as a function of its energy, frequency, or wavelength.
The maximum possible value for the speed of light is that of light in a vacuum, and this speed is used for a constant in many
area of physics.
c is the symbol used to represent the speed of light in a vacuum, and its value is 299,792,458 meters per second.
When light travels through medium, its speed is hindered by the index of refraction of that medium. Its actual speed can be
found with: v = .n
c
Key Terms
differential equation: An equation involving the derivatives of a function.
flux: A quantitative description of the transfer of a given vector quantity through a surface. In this context, we refer to the
electric flux and magnetic flux.
electromagnetic wave: A wave of oscillating electric and magnetic fields.
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phase: Waves are said to be “in phase” when they begin at the same part (e.g., crest) of their respective cycles.
photon: The quantum of light and other electromagnetic energy, regarded as a discrete particle having zero rest mass, no
electric charge, and an indefinitely long lifetime.
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ, and corresponds to the velocity of the wave divided by its frequency.
frequency: The quotient of the number of times n a periodic phenomenon occurs over the time t in which it occurs: f = n / t.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
refractive index: The ratio of the speed of light in air or vacuum to that in another medium.
doppler effect: Apparent change in frequency of a wave when the observer and the source of the wave move relative to each
other.
classical electrodynamics: A branch of theoretical physics that studies consequences of the electromagnetic forces between
electric charges and currents.
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23.3: Applications of EM Waves
learning objectives
Discuss wireless technologies and their applications
Wireless communication is the transfer of information between two or more points that are not connected by an electrical
conductor. The term is commonly used in the telecommunications industry to refer to telecommunications systems (e.g., radio
transmitters and receivers, remote controls, etc.) that use some form of energy (e.g., radio waves, acoustic energy, etc.) to transfer
information without the use of wires. Information is transferred in this manner over both short and long distances. Wireless
operations permit services, such as long-range communications, that are otherwise impossible (or impractical) to implement with
the use of wires.
The most common wireless technologies use electromagnetic wireless telecommunications, such as radio or infra-red signals. With
infra-red waves, distances are short (such as a few meters for television remote control) while radio waves can reach as far as
thousands or even millions of kilometers for deep-space radio communications. It encompasses various types of fixed, mobile, and
portable applications, including two-way radios, cellular telephones, personal digital assistants (PDAs), and wireless networking.
Other examples of applications of radio wireless technology include GPS units, garage door openers, wireless computer mice,
keyboards and headsets, headphones, radio receivers, satellite television, broadcast television, and cordless telephones. Less
common methods of achieving wireless communications include the use of light, sound, magnetic, or electric fields.
Two cellular phones: The Qualcomm QCP-2700, a mid-1990s candybar style phone, and an iPhone 4S, a current production
smartphone.
One of the best-known examples of wireless technology is the mobile (or cellular) phone, with more than 4.6 billion mobile cellular
subscriptions worldwide as of the end of 2010 (examples of such phones are shown in ). These wireless devices use radio waves to
enable their users to make phone calls from many locations worldwide. They can be used within range of the mobile telephone sites
that house the necessary equipment to transmit and receive the radio signals these devices emit. Wireless data communications are
also an essential component of mobile computing. The various available technologies differ in local availability, coverage range,
and performance. In some circumstances, users must be able to employ multiple connection types and switch between them.
To simplify the experience for the user, connection manager software is available, or a mobile VPN can be utilized to handle the
multiple connections as a secure, single virtual network. One popular supporting technology is Wi-Fi, a wireless local area network
that enables portable computing devices to connect easily to the Internet. Standardized as IEEE 802.11 a,b,g,n, Wi-Fi approaches
speeds of some types of wired Ethernet. Wi-Fi has become the de facto standard for access in private homes, within offices, and at
public hotspots. Some businesses charge customers a monthly fee for the service, while others offer it for free in an effort to
increase sales of their goods.
Cellular data service offers coverage within a range of 10-15 miles from the nearest cell site. Speeds have increased as technologies
have evolved, from earlier technologies such as GSM, CDMA and GPRS, to 3G networks such as W-CDMA, EDGE or
CDMA2000. Mobile Satellite Communications may be used where other wireless connections are unavailable, such as in largely
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rural areas or remote locations. Satellite communications are especially important for transportation, aviation, maritime, and
military use.
Key Points
Wireless operations permit services, such as long-range communications, that are impossible or impractical to implement with
the use of wires.
The most common wireless technologies use electromagnetic wireless telecommunications, such as radio.
Less common methods of achieving wireless communications include the use of light, sound, magnetic, or electric fields.
Key Terms
radio wave: Electromagnetic radiation having a wavelength between about.5 centimeters and 30,000 meters; used for the
broadcasting of radio and television signals.
conductor: A material which contains movable electric charges.
telecommunication: The science and technology of the communication or messages over a distance, especially using electric,
electronic or electromagnetic impulses.
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CHAPTER OVERVIEW
24: Geometric Optics is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
24.1: Overview
learning objectives
Distinguish three ways that rays can travel
Rays, or beams of light, can travel in three ways: directly, through a material, or indirectly ( reflection ). These three methods of
light travel are shown in this image. The word ray comes from mathematics, and refers to a straight line that originates at some
point. Even when passing through a material, or bouncing off of a material in a reflection, the light continues to travel in a straight
line, even if that line has changed direction. The movement of light, as a ray, can be shown with simple geometry and trigonometry.
This is called geometric optics.
Methods of Travel by a Ray of Light: Light can travel through empty space directly from the source, through media like air and
glass, reflect from an object like a mirror, or travel in a straight line.
Refraction of Light Rays: The concept of refraction explains how a pencil submerged in water appears to bend.
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Light Bouncing Off a Material
When light is bounced off of a material, such as a mirror, this is called a reflection. This is when a light ray, the incident ray, hits a
reflective material and bounces off as the reflected ray at a specific angle. This is called the angle of reflection. Since the movement
of the light rays can be shown geometrically, if a mirror is one-half your height, you could see your whole body in the reflection.
Reflected Rays: This diagram shows how light rays reflects off of a surface.
Key Points
Direct light motion is when a ray of light travels from a source and is uninterrupted until it reaches it destination.
Refracted light is when a light ray travels through another medium and, due to the difference in refractive indexes of the
materials, changes direction slightly.
When a light ray hits a reflective material, it bounces off as a reflected ray at a specific angle.
Key Terms
refraction: Changing of a light ray’s direction when it passes through variations in matter.
geometric optics: Optics that describes light propagation in terms of “rays”.
reflection: the property of a propagated wave being thrown back from a surface (such as a mirror)
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24.1.2 https://phys.libretexts.org/@go/page/15715
24.2: Reflection, Refraction, and Dispersion
learning objectives
Formulate the relationship between the angle of reflection and the angle of incidence
Whenever you look into a mirror or squint at sunlight glinting off a lake, you are seeing a reflection. When you look at the text in a
book, you are actually seeing the light that is reflected from it. Large telescopes use reflections to form images of stars and other
astronomical objects. In fact, the only way we can see an object that does not itself emit light is if that object reflects light.
The law of reflection is illustrated in, which also shows how the angles are measured relative to the perpendicular to the surface at
the point where the light ray strikes. The law of reflection is very simple: The angle of reflection equals the angle of incidence.
When we see our reflection in a mirror, it appears that our image is actually behind the mirror — we see the light coming from a
direction determined by the law of reflection. The angles are such that our image appears exactly the same distance behind the
mirror as we stand away from the mirror.
Mirror Reflection: An image in a mirror appears as though it is behind the mirror. The two rays shown are those that strike the
mirror at just the correct angles to be reflected into the eyes of the viewer. The image appears to come from the direction the rays
are coming from when they enter the viewer’s eyes.
Law of Reflection: The law of reflection states that the angle of reflection equals the angle of incidence: θr = θi. The angles are
measured relative to the perpendicular to the surface at the point where the ray strikes the surface.
24.2.1 https://phys.libretexts.org/@go/page/15716
We expect to see reflections off a smooth surface. However, light strikes different parts of a rough surface at different angles, and it
is reflected in many different directions (“diffused”). Diffused light is what allows us to see a sheet of paper from any angle. Many
objects, such as people, clothing, leaves, and walls, have rough surfaces and can be seen from all sides. A mirror, on the other hand,
has a smooth surface (compared with the wavelength of light) and reflects light at specific angles. When the moon reflects off the
surface of a lake, a combination of these effects takes place.
learning objectives
Formulate the relationship between the index of refraction and the speed of light
It is easy to notice some odd things when looking into a fish tank. For example, you may see the same fish appearing to be in two
different places. This is because light coming from the fish to us changes direction when it leaves the tank, and in this case, it can
travel two different paths to get to our eyes. The changing of a light ray’s direction (loosely called bending) when it passes through
24.2.2 https://phys.libretexts.org/@go/page/15716
variations in matter is called refraction. Refraction is responsible for a tremendous range of optical phenomena, from the action of
lenses to voice transmission through optical fibers.
Law of Refraction: Looking at the fish tank as shown, we can see the same fish in two different locations, because light changes
directions when it passes from water to air. In this case, the light can reach the observer by two different paths, and so the fish
seems to be in two different places. This bending of light is called refraction and is responsible for many optical phenomena.
Refraction: The changing of a light ray’s direction (loosely called bending) when it passes through variations in matter is called
refraction.
Speed of Light
The speed of light c not only affects refraction, it is one of the central concepts of Einstein’s theory of relativity. The speed of light
varies in a precise manner with the material it traverses. It makes connections between space and time and alters our expectations
that all observers measure the same time for the same event, for example. The speed of light is so important that its value in a
vacuum is one of the most fundamental constants in nature as well as being one of the four fundamental SI units.
Why does light change direction when passing from one material ( medium ) to another? It is because light changes speed when
going from one material to another.
Law of Refraction
A ray of light changes direction when it passes from one medium to another. As before, the angles are measured relative to a
perpendicular to the surface at the point where the light ray crosses it. The change in direction of the light ray depends on how the
speed of light changes. The change in the speed of light is related to the indices of refraction of the media involved. In mediums
that have a greater index of refraction the speed of light is less. Imagine moving your hand through the air and then moving it
through a body of water. It is more difficult to move your hand through the water, and thus your hand slows down if you are
applying the same amount of force. Similarly, light travels slower when moving through mediums that have higher indices of
refraction.
The amount that a light ray changes its direction depends both on the incident angle and the amount that the speed changes. For a
ray at a given incident angle, a large change in speed causes a large change in direction, and thus a large change in angle. The exact
mathematical relationship is the law of refraction, or “Snell’s Law,” which is stated in equation form as:
24.2.3 https://phys.libretexts.org/@go/page/15716
Here n1 and n2 are the indices of refraction for medium 1 and 2, and θ1 and θ2are the angles between the rays and the perpendicular
in medium 1 and 2. The incoming ray is called the incident ray and the outgoing ray the refracted ray, and the associated angles the
incident angle and the refracted angle. The law of refraction is also called Snell’s law after the Dutch mathematician Willebrord
Snell, who discovered it in 1621. Snell’s experiments showed that the law of refraction was obeyed and that a characteristic index
of refraction n could be assigned to a given medium.
Understanding Snell’s Law with the Index of Refraction: This video introduces refraction with Snell’s Law and the index of
refraction.The second video discusses total internal reflection (TIR) in detail. http://www.youtube.com/watch?v=fvrvqm3Erzk
learning objectives
Formulate conditions required for the total internal reflection
24.2.4 https://phys.libretexts.org/@go/page/15716
Total internal reflection is a phenomenon that happens when a propagating wave strikes a medium boundary at an angle larger than
a particular critical angle with respect to the normal to the surface. If the refractive index is lower on the other side of the boundary
and the incident angle is greater than the critical angle, the wave cannot pass through and is entirely reflected. The critical angle is
the angle of incidence above which the total internal reflectance occurs.
What is Total Internal Reflection?: Describes the concept of total internal reflection, derives the equation for the critical angle
and shows one example.
Critical angle
The critical angle is the angle of incidence above which total internal reflection occurs. The angle of incidence is measured with
respect to the normal at the refractive boundary (see diagram illustrating Snell’s law ). Consider a light ray passing from glass into
air. The light emanating from the interface is bent towards the glass. When the incident angle is increased sufficiently, the
transmitted angle (in air) reaches 90 degrees. It is at this point no light is transmitted into air. The critical angle θcθc is given by
Snell’s law, n1sinθ1=n2sinθ2n1sinθ1=n2sinθ2. Here, n1 and n2 are refractive indices of the media, and θ1θ1 and θ2θ2are angles of
incidence and refraction, respectively. To find the critical angle, we find the value for θ1θ1 when θ2θ2= 90° and thus
sinθ2=1sinθ2=1. The resulting value of θ1θ1 is equal to the critical angle θc=θ1=arcsin(n2n1)θc=θ1=arcsin(n2n1). So the critical
angle is only defined when n2/n1 is less than 1.
24.2.5 https://phys.libretexts.org/@go/page/15716
Fig 1: Refraction of light at the interface between two media, including total internal reflection.
Optical Fiber
Total internal reflection is a powerful tool since it can be used to confine light. One of the most common applications of total
internal reflection is in fibre optics. An optical fibre is a thin, transparent fibre, usually made of glass or plastic, for transmitting
light. The construction of a single optical fibre is shown in.
Fig 2: Fibers in bundles are clad by a material that has a lower index of refraction than the core to ensure total internal reflection,
even when fibers are in contact with one another. This shows a single fiber with its cladding.
The basic functional structure of an optical fiber consists of an outer protective cladding and an inner core through which light
pulses travel. The overall diameter of the fiber is about 125 μm and that of the core is just about 50 μm. The difference in refractive
index of the cladding and the core allows total internal reflection in the same way as happens at an air-water surface show in. If
light is incident on a cable end with an angle of incidence greater than the critical angle then the light will remain trapped inside the
glass strand. In this way, light travels very quickly down the length of the cable over a very long distance (tens of kilometers).
Optical fibers are commonly used in telecommunications, because information can be transported over long distances, with
minimal loss of data. Another common use can be found in medicine in endoscopes. The field of applied science and engineering
concerned with the design and application of optical fibers are called fiber optics.
Total Polarization
Brewster’s angle is an angle of incidence at which light with a particular polarization is perfectly transmitted through a surface.
learning objectives
Calculate the Brewster’s angle from the indices of refraction and discuss its physical mechanism
24.2.6 https://phys.libretexts.org/@go/page/15716
Brewster’s angle (also known as the polarization angle) is an angle of incidence at which light with a particular polarization is
perfectly transmitted through a transparent dielectric surface, with no reflection. When unpolarized light is incident at this angle,
the light that is reflected from the surface is therefore perfectly polarized. This special angle of incidence is named after the
Scottish physicist Sir David Brewster (1781–1868).
The physical mechanism for this can be qualitatively understood from the manner in which electric dipoles in the media respond to
p-polarized light (whose electric field is polarized in the same plane as the incident ray and the surface normal). One can imagine
that light incident on the surface is absorbed, and then re-radiated by oscillating electric dipoles at the interface between the two
media. The refracted light is emitted perpendicular to the direction of the dipole moment; no energy can be radiated in the direction
of the dipole moment. Thus, if the angle of reflection θ1 (angle of reflection) is equal to the alignment of the dipoles (90 – θ2),
where θ2 is angle of refraction, no light is reflected.
Fig 1: An illustration of the polarization of light that is incident on an interface at Brewster’s angle.
This geometric condition can be expressed as θ + θ = 90 , where θ1 is the angle of incidence and θ2 is the angle of refraction.
1 2
∘
Using Snell’s law (n1sinθ1 = n2sinθ2), one can calculate the incident angle θ1 = B at which no light is reflected:
n2
n1 sin(θB ) = n2 sin(90
∘
− θB ) Solving for θB gives θ B = arctan(
n1
) .
When light hits a surface at a Brewster angle, reflected beam is linearly polarized. shows an example, where the reflected beam was
nearly perfectly polarized and hence, blocked by a polarizer on the right picture. Polarized sunglasses use the same principle to
reduce glare from the sun reflecting off horizontal surfaces such as water or road.
Fig 2: Photograph taken of a window with a camera polarizer filter rotated to two different angles. In the picture at left, the
polarizer is aligned with the polarization angle of the window reflection. In the picture at right, the polarizer has been rotated 90°
eliminating the heavily polarized reflected sunlight.
24.2.7 https://phys.libretexts.org/@go/page/15716
Joy Physics- exp #26 Polarization
Polarization Experience: A polarizing filter allows light of a particular plane of polarization to pass, but scatters the rest of the
light. When two polarizing filters are crossed, almost no light gets through. Some materials have molecules that rotate the plane of
polarization of light. When one of these materials is placed between crossed polarizing filters, more light is allowed to pass
through.
learning objectives
Describe production of rainbows by a combination of refraction and reflection processes
We see about six colors in a rainbow—red, orange, yellow, green, blue, and violet; sometimes indigo is listed, too. These colors are
associated with different wavelengths of light. White light, in particular, is a fairly uniform mixture of all visible wavelengths.
Sunlight, considered to be white, actually appears to be a bit yellow because of its mixture of wavelengths, but it does contain all
visible wavelengths. The sequence of colors in rainbows is the same sequence as the colors plotted versus wavelength. What this
implies is that white light is spread out according to wavelength in a rainbow. Dispersion is defined as the spreading of white light
into its full spectrum of wavelengths. More technically, dispersion occurs whenever there is a process that changes the direction of
light in a manner that depends on wavelength. Dispersion, as a general phenomenon, can occur for any type of wave and always
involves wavelength-dependent processes.
24.2.8 https://phys.libretexts.org/@go/page/15716
Colors of a Rainbow: Even though rainbows are associated with seven colors, the rainbow is a continuous distribution of colors
according to wavelengths.
Refraction is responsible for dispersion in rainbows and many other situations. The angle of refraction depends on the index of
refraction, as we saw in the Law of Refraction. We know that the index of refraction n depends on the medium. But for a given
medium, n also depends on wavelength. Note that, for a given medium, n increases as wavelength decreases and is greatest for
violet light. Thus violet light is bent more than red light and the light is dispersed into the same sequence of wavelengths.
Pure Light and Light Dispersion: (a) A pure wavelength of light falls onto a prism and is refracted at both surfaces. (b) White
light is dispersed by the prism (shown exaggerated). Since the index of refraction varies with wavelength, the angles of refraction
vary with wavelength. A sequence of red to violet is produced, because the index of refraction increases steadily with decreasing
wavelength.
Rainbows are produced by a combination of refraction and reflection. You may have noticed that you see a rainbow only when you
look away from the sun. Light enters a drop of water and is reflected from the back of the drop. The light is refracted both as it
enters and as it leaves the drop. Since the index of refraction of water varies with wavelength, the light is dispersed, and a rainbow
is observed. (There is no dispersion caused by reflection at the back surface, since the law of reflection does not depend on
wavelength. ) The actual rainbow of colors seen by an observer depends on the myriad of rays being refracted and reflected toward
24.2.9 https://phys.libretexts.org/@go/page/15716
the observer’s eyes from numerous drops of water. The arc of a rainbow comes from the need to be looking at a specific angle
relative to the direction of the sun.
Light Reflecting on Water Droplet: Part of the light falling on this water drop enters and is reflected from the back of the drop.
This light is refracted and dispersed both as it enters and as it leaves the drop.
Key Points
Light strikes different parts of a rough surface at different angles and is reflected, or diffused, in many different directions.
A mirror has a smooth surface (compared with the wavelength of light) and so reflects light at specific angles.
We see the light reflected off a mirror coming from a direction determined by the law of reflection.
The changing of a light ray’s direction (loosely called bending) when it passes through variations in matter is called refraction.
The index of refraction is n=c/v, where v is the speed of light in the material, c is the speed of light in vacuum, and n is the
index of refraction.
Snell’s law, the law of refraction, is stated in equation form as n1sinθ1=n2sinθ2n1sinθ1=n2sinθ2.
The critical angle is the angle of incidence above which total internal reflection occurs and given as
θc=arcsin(n2n1)θc=arcsin(n2n1).
The critical angle is only defined when n2/n1 is less than 1.
If light is incident on an optical fiber with an angle of incidence greater than the critical angle then the light will remain trapped
inside the glass strand. Light can travel over a very long distance without a significant loss.
When light hits a surface at a Brewster angle, reflected beam is linearly polarized.
The physical mechanism for the Brewster’s angle can be qualitatively understood from the manner in which electric dipoles in
the media respond to p-polarized light.
Brewsters’ angle is given as θB=arctan(n2n1)θB=arctan(n2n1).
Dispersion occurs whenever there is a process that changes the direction of light in a manner that depends on wavelength.
Dispersion can occur for any type of wave and always involves wavelength-dependent processes.
For a given medium, n increases as wavelength decreases and is greatest for violet light. Thus violet light is bent more than red
light, as can be seen with a prism.
In a rainbow, light enters a drop of water and is reflected from the back of the drop. The light is refracted both as it enters and as
it leaves the drop.
Key Terms
reflection: the property of a propagated wave being thrown back from a surface (such as a mirror)
refraction: Changing of a light ray’s direction when it passes through variations in matter.
index of refraction: For a material, the ratio of the speed of light in vacuum to that in the material.
Snell’s law: A formula used to describe the relationship between the angles of incidence and refraction.
cladding: One or more layers of materials of lower refractive index, in intimate contact with a core material of higher refractive
index.
dipole: A separation of positive and negative charges.
dielectric: An electrically insulating or nonconducting material considered for its electric susceptibility (i.e., its property of
polarization when exposed to an external electric field).
polarizer: An optical filter that passes light of a specific polarization and blocks waves of other polarizations.
24.2.10 https://phys.libretexts.org/@go/page/15716
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24.3: Lenses
learning objectives
Describe properties of a thin lens and the purpose of ray tracing
Convex Lens: Rays of light entering a converging lens parallel to its axis converge at its focal point F. (Ray 2 lies on the axis of the
lens. ) The distance from the center of the lens to the focal point is the lens’s focal length f. An expanded view of the path taken by
ray 1 shows the perpendiculars and the angles of incidence and refraction at both surfaces.
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Thin Lens: Thin lenses have the same focal length on either side. (a) Parallel light rays entering a converging lens from the right
cross at its focal point on the left. (b) Parallel light rays entering a diverging lens from the right seem to come from the focal point
on the right.
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Diverging Lens: Rays of light entering a diverging lens parallel to its axis are diverged, and all appear to originate at its focal point
F. The dashed lines are not rays—they indicate the directions from which the rays appear to come. The focal length f of a diverging
lens is negative. An expanded view of the path taken by ray 1 shows the perpendiculars and the angles of incidence and refraction
at both surfaces.
24.3.3 https://phys.libretexts.org/@go/page/15717
Ray Diagrams, Concave Lens and Convex Mirror: Shows how to draw the ray diagrams for locating the image produced by a
concave lens and a convex mirror.
learning objectives
Formulate five basic rules of ray tracing
24.3.4 https://phys.libretexts.org/@go/page/15717
Image Formation with a Thin Lens: Ray tracing is used to locate the image formed by a lens. Rays originating from the same
point on the object are traced—the three chosen rays each follow one of the rules for ray tracing, so that their paths are easy to
determine. The image is located at the point where the rays cross. In this case, a real image—one that can be projected on a screen
—is formed.
Several important distances appear in. We define do as the object distance—the distance of an object from the center of a lens.
Image distance diis defined as the distance of the image from the center of a lens. The height of the object and height of the image
are given the symbols ho and hi, respectively. Images that appear upright relative to the object have heights that are positive and
those that are inverted have negative heights. Using the rules of ray tracing and making a scale drawing with paper and pencil, like
that in, we can accurately describe the location and size of an image. But the real benefit of ray tracing is in visualizing how images
are formed in a variety of situations. To obtain numerical information, we use a pair of equations that can be derived from a
geometric analysis of ray tracing for thin lenses. The thin lens equation is:
1 1 1
+ = (24.3.1)
do di f
We define the ratio of image height to object height (hi/ho) as the magnification m. The magnification is related to do, di, ho, and hi
by the following relation:
hi di
=− =m (24.3.2)
ho do
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In many cases both of these equations are referred to together as the thin lens equations. The thin lens equations are broadly
applicable to all situations involving thin lenses (and “thin” mirrors).
Thin Lens Equations for a Convex Lens: Shows how to use the thin lens equation to calculate the image distance, image height
and image orientation for convex lenses when the object distance is greater the the focal length (f).
Combinations of Lenses
A compound lens is an array of simple lenses with a common axis.
learning objectives
Calculate focal length for a compound lens from the focal lengths of simple lenses
COMPOUND LENSES
In contrast to a simple lens, which consists of only one optical element, a compound lens is an array of simple lenses (elements)
with a common axis. The use of multiple elements allows for the correction of more optical aberrations, such as the chromatic
aberration caused by the wavelength-dependent index of refraction in glass, than is possible using a single lens. In many cases these
aberrations can be compensated for to a great extent by using a combination of simple lenses with complementary aberrations.
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The simplest case is where lenses are placed in contact: if the lenses of focal lengths f1 and f2 are “thin”, the combined focal length
f of the lenses is given by
1 1 1
= + (24.3.3)
f f1 f2
Since 1/f is the power of a lens, it can be seen that the powers of thin lenses in contact are additive.
If two thin lenses are separated in air by some distance d (where d is smaller than the focal length of the first lens), the focal length
for the combined system is given by
1 1 1 d
= + − (24.3.4)
f f1 f2 f1 f2
As d tends to zero, the value of the BFL tends to the value of f given for thin lenses in contact.
If the separation distance is equal to the sum of the focal lengths (d = f1+f2), the combined focal length and BFL are infinite. This
corresponds to a pair of lenses that transform a parallel (collimated) beam into another collimated beam (see ). This type of system
is called an afocal system, since it produces no net convergence or divergence of the beam. Two lenses at this separation form the
simplest type of optical telescope. Although the system does not alter the divergence of a collimated beam, it does alter the width of
the beam. The magnification of such a telescope is given by
Keplerian Telescope: All refracting telescopes use the same principles. The combination of an objective lens 1 and some type of
eyepiece 2 is used to gather more light than the human eye could collect on its own, focus it 5, and present the viewer with a
brighter, clearer, and magnified virtual image 6. The magnification can be found by dividing the focal length of the objective lens
by the focal length of the eyepiece.
f2
M =− (24.3.6)
f1
which is the ratio of the input beam width to the output beam width. Note the sign convention: a telescope with two convex lenses
(f1 > 0, f2 > 0) produces a negative magnification, indicating an inverted image. A convex plus a concave lens (f1 > 0 >f2) produces
a positive magnification and the image is upright.
ACHROMATS
An achromatic lens or achromat is a lens that is designed to limit the effects of chromatic and spherical aberration. Achromatic
lenses are corrected to bring two wavelengths (typically red and blue/violet) into focus in the same plane.
The most common type of achromat is the achromatic doublet, which is composed of two individual lenses made from glasses with
different amounts of dispersion Typically, one element is a negative (concave) element made out of flint, which has relatively high
dispersion, and the other is a positive (convex) element made of crown glass, which has lower dispersion. The lens elements are
mounted next to each other, often cemented together, and shaped so that the chromatic aberration of one is counterbalanced by that
of the other.
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In the most common type (shown in ), the positive power of the crown lens element is not quite equaled by the negative power of
the flint lens element. Together they form a weak positive lens that will bring two different wavelengths of light to a common
focus. Negative doublets, in which the negative-power element predominates, are also made.
Achromatic Doublet: (a) Chromatic aberration is caused by the dependence of a lens’s index of refraction on color (wavelength).
The lens is more powerful for violet (V) than for red (R), producing images with different locations and magnifications. (b)
Multiple-lens systems, such as this achromatic doublet, can partially correct chromatic aberrations, but they may require lenses of
different materials and add to the expense of optical systems such as cameras.
learning objectives
Compare idealized thin lenses with real lenses
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Thick Converging Lens: Diagram of a positive (converging) lens. The lensmaker’s formula relates the radii of curvature, the
index of refraction of the lens, the thickness of the lens, and the focal length.
If the lens is biconcave, a beam of light passing through the lens is diverged (spread); the lens is thus called a negative or diverging
lens. The beam after passing through the lens appears to be emanating from a particular point on the axis in front of the lens; the
distance from this point to the lens is also known as the focal length, although it is negative with respect to the focal length of a
converging lens. See for a diagram of a negative (diverging) lens.
Negative Diverging Lens: Diagram of a negative (diverging) lens. The lensmaker’s formula relates the radii of curvature, the
index of refraction of the lens, the thickness of the lens, and the focal length.
The focal length of a thick lens in air can be calculated from the lensmaker’s equation:
1 1 1 (n − 1)d
P = = (n − 1) [ − + ] (24.3.7)
f R1 R2 nR 1 R2
where
P is the power of the lens,
f is the focal length of the lens,
n is the refractive index of the lens material,
R1is the radius of curvature of the lens surface closest to the light source,
R2 is the radius of curvature of the lens surface farthest from the light source, d and is the thickness of the lens (the distance
along the lens axis between the two surface vertices).
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Thin Lens Approximation
The above equation can be greatly simplified if the lens thickness d is very small compared to R1 and R2. In this case, the thin lens
approximation can then be made and the lensmaker’s equation can be approximated as
1 1 1
P = ≈ (n − 1) [ − ] (24.3.8)
f R1 R2
The focal length f is positive for converging lenses, and negative for diverging lenses. The reciprocal of the focal length, 1/f, is the
optical of the lens. If the focal length is in meters, this gives the optical power in diopters (inverse meters).
Lenses have the same focal length when light travels from the back to the front as when light goes from the front to the back,
although other properties of the lens, such as the aberrations are not necessarily the same in both directions.
learning objectives
Compare the effect of a convex lens and a concave lens on the light rays
Magnifying Glass: Sunlight focused by a converging magnifying glass can burn paper. Light rays from the sun are nearly parallel
and cross at the focal point of the lens. The more powerful the lens, the closer to the lens the rays will cross.
The greater effect a lens has on light rays, the more powerful it is said to be. For example, a powerful converging lens will focus
parallel light rays closer to itself and will have a smaller focal length than a weak lens. The light will also focus into a smaller, more
24.3.10 https://phys.libretexts.org/@go/page/15717
intense spot for a more powerful lens. The power P of a lens is defined as the inverse of its focal length. In equation form:
1
P = (24.3.9)
f
shows the effect of a concave lens on rays of light entering it parallel to its axis (the path taken by ray 2 in the figure is the axis of
the lens). The concave lens is a diverging lens, because it causes the light rays to bend away (diverge) from its axis. In this case, the
lens is shaped so that all light rays entering it parallel to its axis appear to originate from the same point F, defined as the focal
point of a diverging lens. The distance from the center of the lens to the focal point is again called the focal length f of the lens.
Note that the focal length and power of a diverging lens are defined as negative. For example, if the distance to F in is 5.00 cm,
then the focal length is f=–5.00 cm and the power of the lens is P=–20 D. The expanded view of the path of one ray through the
lens illustrates how the shape of the lens (given the law of refraction) causes the ray to follow its particular path and be diverged.
In subsequent sections we will examine the technique of ray tracing to describe the formation of images by lenses. Additionally, we
will explore how image locations and characteristics can be quantified with the help of a set of geometric optics equations.
Key Points
When light interacts with objects several times as large as its wavelength, it travels in straight lines and acts like a ray. A ray is
simply a straight line that originates at a point.
Ray tracing is the method for determining the paths light takes through matter, such as optical systems that include lenses.
A thin lens is defined as one with a thickness that allows rays to refract, as illustrated in, but that does not allow properties such
as dispersion and aberrations. An ideal thin lens has two refracting surfaces but the lens is thin enough to assume that light rays
bend only once.
There are five basic rules for tracing rays through a lens.
Ray tracing can be used to construct an image from the light rays originating from an object that pass through a lens. The image
is located at the point where the rays cross. By choosing several points from an object the entire image can be constructed.
We define do to be the object distance, the distance of an object from the center of a lens. Image distance di is defined to be the
distance of the image from the center of a lens. The height of the object and height of the image are given the symbols ho and
hi, respectively.
The thin lens equation quickly provides the relation between di, do, and the focal length f. It can be derived from a geometric
analysis of ray tracing for thin lenses and is given by 1
do
+
1
=
di
1
f
.
The magnification m of an image is the ratio between the image and object height (hi/ho). The magnification is related to do, di,
hi di
ho, and hi by the following relation: ho
=−
do
=m .
The use of multiple elements allows for the correction of more optical aberrations, such as the chromatic abberation caused by
the wavelength -dependent index of refraction in glass, than is possible using a single lens.
If the lenses of focal lengths f1 and f2are “thin”, the combined focal length f of the lenses is given by = + while if the
1
f
1
f1
1
f2
lenses are separated by some distance d then the combined focal length is given by 1
f
=
1
f
+
1
f
−
d
f f
.
1 2 1 2
If the separation distance is equal to the sum of the focal lengths (d = f1+f2), the combined focal length is infinite. This
corresponds to a pair of lenses that transform a collimated beam into another collimated beam. This type of system is called an
afocal system (a simple optical telescope).
An achromatic doublet is a kind of compound lens designed to bring two wavelengths (typically red and blue/violet) into focus
in the same plane. This (partially) corrects for the chromatic aberration found in a single simple lens. See.
If a lens is biconvex, a beam of light travelling parallel to the lens axis and passing through the lens will be focused to a spot on
the axis, at a certain distance behind the lens (i.e. the focal length ). In this case, the lens is called a positive or converging lens.
See.
If a lens is biconcave, a beam of light passing through the lens is diverged (spread); the lens is thus called a negative or
diverging lens. See.
(n−1)d
The focal length of a thick lens in air can be calculated from the lensmaker’s equation: P = 1
f
= (n − 1) [
1
R1
−
1
R2
+
nR 1 R2
] .
The signs of the lens’ radii of curvature indicate whether the corresponding surfaces are convex or concave. The signs are
reversed for the back surface of the lens: if R2 is positive the surface is concave, and if R2 is negative the surface is convex.
The lensmaker’s equation can be greatly simplified if the lens thickness d is very small compared to R1 and R2. In this case, the
thin lens approximation can then be made and the lensmaker’s equation can be approximated as P =
1
f
≈ (n − 1) [
1
R1
−
1
R2
] .
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Recall that the a ray will bend as it enters a medium with a different refractive index. Since the refractive index of a lens is
greater than air, a light ray will move towards the perpendicular as it enters and away as it leaves.
A convex lens has been shaped so that all light rays that enter it parallel to its axis cross one another at a single point on the
opposite side of the lens (the focal point ). Such a lens is called a converging (or convex) lens for the converging effect it has on
light rays. See.
A concave lens is a diverging lens, because it causes the light rays to bend away (diverge) from its axis. shows the effect it has
on rays of light that enter it parallel to its axis (the path taken by ray 2 in the figure is the axis of the lens).
The greater effect a lens has on light rays, the more powerful it is said to be. A powerful converging lens will focus parallel
light rays closer to itself and will have a smaller focal length than a weak lens. The power of a lens is given by the equation
P =
1
f
.
Key Terms
focal point: A focus—a point at which rays of light or other radiation converge.
ray tracing: A technique used in optics for analysis of optical systems.
thin lens: A thin lens is defined to be one whose thickness allows rays to refract but does not allow properties such as
dispersion and aberrations.
thin lens equation: Relates object distance do, image distance di, and focal length f: 1
+
do
1
di
=
1
image distance: The distance of the image from the center of the lens.
magnification: The apparent enlargement of an object in an image.
aberration: The convergence to different foci, by a lens or mirror, of rays of light emanating from one and the same point, or
the deviation of such rays from a single focus; a defect in a focusing mechanism that prevents the intended focal point.
afocal system: An optical system that produces no net convergence or divergence of the beam, i.e. has an infinite effective focal
length. This type of system can be created with a pair of optical elements where the distance between the elements is equal to
the sum of each element’s focal length (d = f1+f2).
achromatic doublet: A type of lens made up of two simple lenses paired together designed so that the chromatic aberration of
each lens partially offsets the other; in this way light in a range of wavelengths may be brought to the same focus.
thick lens: Lenses whose thicknesses are not negligible (i.e., one cannot make the simple assumption that a light ray is refracted
only once in the lens).
surface vertices: The points where each surface crosses the optical axis. They are important primarily because they are the
physically measurable parameters for the position of the optical elements, and so the positions of the other cardinal points must
be known with respect to the vertices to describe the physical system.
convex lens: A lens having at least one convex surface, such that light passing through it, may be brought to a focus.
concave lens: A lens having at least one concave surface, such that light rays passing through it bend away from its optical axis.
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24.4: Mirrors
learning objectives
Describe interaction of the light with a mirror surface
Ray Diagrams
The way that we can predict how a reflection will look is by drawing a ray diagram. These diagrams can be used to find the
position and size of the image and whether that image is real or virtual. These are the steps you follow to draw a ray diagram:
1. Draw the plane mirror as a straight line on a principal axis. The principal axis is an imaginary line that is drawn perpendicular
to the mirror.
2. Draw the object as an arrow in front of the mirror.
3. Draw the image of the object, by using the principle that the image is placed at the same distance behind the mirror that the
object is in front of the mirror. The image size is also the same as the object size. shows these first three steps.
4. Place a dot at the point the eye is located.
5. Pick one point on the image and draw the reflected ray that travels to the eye as it sees this point. Remember to add an
arrowhead.
6. Draw the incident ray for light traveling from the corresponding point on the object to the mirror, such that the law of reflection
is obeyed.
7. Continue for other extreme points on the object (i.e. the tip and base of the arrow). A completed ray diagram is shown in
The angle in which a light ray hits the mirror is the same angle in which it will be reflected back. If, for example, a light ray leaves
the top of an object travelling parallel to the principal axis, it will hit the mirror at a 0 degree angle, and be reflected back at 0
degrees. When this happens, we say the ray hit the mirror normally. If the light ray hit the object at a 30 degree angle, it will be
reflected back at a 30 degree angle.
learning objectives
Distinguish properties of the concave and the convex mirrors
Overview
A mirror is a reflective surface that light does not pass through, but bounces off of and this produces an image. Mirrors are made by
putting a thin layer of silver nitrate or aluminium behind a flat piece of glass.
When you place an object in front of a mirror, you see the same object in the mirror. This image that appears to be behind the
mirror is called the image. The object is the source of the incident rays, and the image is formed by the reflected rays. An image
24.4.1 https://phys.libretexts.org/@go/page/16184
formed by reflection may be real or virtual. A real image occurs when light rays actually intersect at the image, and is inverted, or
upside down. A virtual image occurs when light rays do not actually meet at the image. Instead, you “see” the image because your
eye projects light rays backward. A virtual image is right side up (upright).
This section will cover spherical mirrors. Spherical mirrors can be either concave or convex. The center of curvature is the point at
the center of the sphere and describes how big the sphere is. These concepts are shown in.
Spherical Mirrors: This figure shows the difference between a concave and convex mirror.
Concave Mirrors
In a concave mirror, the principal axis is a line that is perpendicular to the center of the mirror. The easiest way to visualize what a
image will look like in this type of mirror is a ray diagram. Before that can be done, the focal point must first be defined. This point
is half way between the mirror and the center of curvature on the principal axis. The distance to the focal point from the mirror is
called the focal length. We can see from the figure that this focal length is also equal to half of the radius of the curvature. shows
the ray diagram of a concave mirror.
Concave Ray Diagram: This is a ray diagram of a concave mirror. The steps taken to draw are the same as those in a plane mirror.
A summary of the properties of concave mirrors is shown below:
converging
real image
inverted
image in front of mirror
Convex Mirrors
In convex mirrors, the principal axis is the same as in a plane or concave mirror, perpendicular to the center of the mirror. In this
case, the focal point is behind the mirror. A convex mirror has a negative focal length because of this. The focal point is the same
distance from the mirror as in a concave mirror. This is shown in.
24.4.2 https://phys.libretexts.org/@go/page/16184
Convex Mirror Ray Diagram: A convex mirror with three rays drawn to locate the image. Each incident ray is reflected
according to the Law of Reflection. The reflected rays diverge. If the reflected rays are extended behind the mirror, then their
intersection gives the location of the image behind the mirror. For a convex mirror, the image is virtual and upright.
A summary of the properties of convex mirrors is shown below:
diverging
virtual image
upright
image behind mirror
Key Points
Reflected images can be either real or virtual. In a plane mirror, the images are virtual.
The virtual images in a plane mirror have a left-right inversion.
Drawing a ray diagram is a way to predict what a reflected image will look like.
Images in mirrors can be either real or virtual.
A summary of the properties of the concave mirrors are shown below: converging real image inverted image in front of mirror.
A summary of the properties of the convex mirrors are shown below: diverging virtual image upright image behind mirror.
Key Terms
virtual image: A virtual image occurs when light rays do not actually meet at the image
concave: curved like the inner surface of a sphere or bowl
convex: curved or bowed outward like the outside of a bowl or sphere or circle
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CHAPTER OVERVIEW
25: Vision and Optical Instruments is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
25.1: The Human Eye
learning objectives
Identify parts of human eye and their functions
The human eye is the gateway to one of our five senses. The human eye is an organ that reacts with light. It allows light perception,
color vision and depth perception. A normal human eye can see about 10 million different colors! There are many parts of a human
eye, and that is what we are going to cover in this atom.
Properties
Contrary to what you might think, the human eye is not a perfect sphere, but is made up of two differently shaped pieces, the
cornea and the sclera. These two parts are connected by a ring called the limbus. The part of the eye that is seen is the iris, which is
the colorful part of the eye. In the middle of the iris is the pupil, the black dot that changes size. The cornea covers these elements,
but is transparent. The fundus is on the opposite of the pupil, but inside the eye and can not be seen without special instruments.
The optic nerve is what conveys the signals of the eye to the brain. is a diagram of the eye. The human eye is made up of three
coats:
Diagram of the Human Eye: The cornea and lens of an eye act together to form a real image on the light-sensing retina, which has
its densest concentration of receptors in the fovea and a blind spot over the optic nerve. The power of the lens of an eye is
adjustable to provide an image on the retina for varying object distances. Layers of tissues with varying indices of refraction in the
lens are shown here. However, they have been omitted from other pictures for clarity.
1. Outermost Layer – composed of the cornea and the sclera.
2. Middle Layer – composed of the choroid, ciliary body and iris.
3. Innermost Layer – the retina, which can be seen with an instrument called the ophthalmoscope.
Once you are inside these three layers, there is the aqueous humor (clear fluid that is contained in the anterior chamber and
posterior chamber), vitreous body (clear jelly that is much bigger than the aqueous humor), and the flexible lens. All of these are
connected by the pupil.
Dynamics
Whenever the eye moves, even just a little, it automatically readjusts the exposure by adjusting the iris, which regulates the size of
the pupil. This is what helps the eye adjust to dark places or really bright lights. The lens of the eye is similar to one in glasses or
cameras. The human eye is had an aperture, just like a camera. The pupil serves this function, and the iris is the aperture stop. The
different parts of the eye has different refractive indexes, and this is what bends the rays to form an image. The cornea provides
two-thirds of the power to the eye. The lens provides the remaining power. The image passes through several layers of the eye, but
happens in a way very similar to that of a convex lens. When the image finally reaches the retena, it is inverted, but the brain will
correct this. shows what happens.
25.1.1 https://phys.libretexts.org/@go/page/16186
Vision Diagram: An image is formed on the retina with light rays converging most at the cornea and upon entering and exiting the
lens. Rays from the top and bottom of the object are traced and produce an inverted real image on the retina. The distance to the
object is drawn smaller than scale.
Eye Movement
Each eye has six muscles; lateral rectus, medial rectus, inferior rectus, superior rectus, inferior oblique, and superior oblique. All of
these muscles provide differnt tensions and torques to control the movement of the eye. These are a few examples of types of eye
movement:
Rapid Eye Movement – Often referred to as REM, this happens in the sleep stage when most vivid dreams occur.
Saccade – These are quick, simultaneous movements of both eyes, and is controlled by the frontal lobe of the brain.
Vestibulo-ocular Reflex – This is the eye movement which is opposite to the movement of the head and keeps the object you are
looking at in the center of vision.
Pursuit Movement – This is the tracking movement when you are following a moving object. It is less accurate than the
vestibulo-ocular reflex.
Color Vision
Using the cone cells in the retina, we perceive images in color; each type of cone specifically sees in regions of red, green, or blue.
learning objectives
Explain how the human eye perceives colors
With human eyesight, cone cells are responsible for color vision. From there, it is important to understand how color is perceived.
Using the cone cells in the retina, we perceive images in color. Each type of cone specifically sees in regions of red, green, or blue,
(RGB), in the color spectrum of red, orange, yellow, green, blue, indigo, violet.
The colors in between these absolutes are seen as different linear combinations of RGB. This is why TVs and computer screens are
made up of thousands of little red, green, or blue lights, and why colors in electronic form are represented by different values of
RGB. These values are usually given in the value of their frequency in log form.
The prior equation deals with the luminance, but the chrominance (dealing with colors) can be found from the following equations:
U = 0.5(BY)
U = 0.5(B − Y) (25.1.2)
V = 0.625(R − Y)
25.1.2 https://phys.libretexts.org/@go/page/16186
V = 0.625(RY) (25.1.3)
You can go from RGB to YUV color spaces with the following matrix operation:
Y R
⎛ ⎞ ⎛ ⎞
⎝ ⎠ ⎝ ⎠
V B
⎝ ⎠
0.4375 −0.3750 −0.0625
Visual Sensitivity
In, we can see that
Visual Sensitivity: This graph shows the sensitivity of the eye to luminance (Y) and chrominance (U, V) components of images.
The horizontal scale is spatial frequency, and represents the frequency of an alternating pattern of parallel stripes with sinusoidally
varying intensity. The vertical scale is the contrast sensitivity of human vision, which is the ratio of the maximum visible range of
intensities to the minimum discernible peak-to-peak intensity variation at the specified frequency.
the maximum sensitivity to Y occurs for spatial frequencies around 5 cycles / degree, which corresponds to striped patterns with
a half-period (stripe width) of 1.8 mm at a distance of 1 m (~arm’s length).
The eye has very little response above 100 cycles / degree, which corresponds to a stripe width of 0.1 mm at 1 m. On a standard
PC display of width 250 mm, this would require 2500 pels per line! Hence the current SVGA standard of 1024×768 pels still
falls somewhat short of the ideal and is limited by CRT spot size. Modern laptop displays have a pel size of about 0.3 mm, but
are pleasing to view because the pel edges are so sharp (and there is no flicker).
The sensitivity to luminance drops off at low spatial frequencies, showing that we are not very good at estimating absolute
luminance levels as long as they do not change with time – the luminance sensitivity to temporal fluctuations (flicker) does not
fall off at low spatial frequencies.
The maximum chrominance sensitivity is much lower than the maximum luminance sensitivity with blue-yellow (U) sensitivity
being about half of red-green (V) sensitivity and about 16 of the maximum luminance sensitivity.
The chrominance sensitivities fall off above 1 cycle / degree, requiring a much lower spatial bandwidth than luminance.
We can now see why it is better to convert to the YUV domain before attempting image compression. The U and V components
may be sampled at a lower rate than Y (due to narrower bandwidth) and may be quantified more coarsely (due to lower contrast
sensitivity).
25.1.3 https://phys.libretexts.org/@go/page/16186
learning objectives
Describe field of view and color sensitivity of the human eye
The human eye is an organ that reacts to light in many circumstances. As a conscious sense organ the human eye allows vision; rod
and cone cells in the retina allow conscious light perception and vision, including color differentiation and the perception of depth.
The human eye can distinguish about 10 million colors. A model of the human eye can be seen in.
Schematic Diagram of the Human Eye: Structure of the eye and closeup of the retina.
The retina of human eye has a static contrast ratio of around 100:1 (about 6.5 f-stops). As soon as the eye moves, it re-adjusts its
exposure, both chemically and geometrically, by adjusting the iris (which regulates the size of the pupil ). Initial dark adaptation
takes place in approximately four seconds of profound, uninterrupted darkness; full adaptation, through adjustments in retinal
chemistry, is mostly complete in thirty minutes. Hence, a dynamic contrast ratio of about 1,000,000:1 (about 20 f-stops) is possible.
The process is nonlinear and multifaceted, so an interruption by light starts the adaptation process over again. Full adaptation is
dependent on good blood flow (thus dark adaptation may be hampered by poor circulation, and vasoconstrictors like tobacco).
The eye includes a lens not dissimilar to lenses found in optical instruments (such as cameras). The same principles can be applied.
The pupil of the human eye is its aperture. The iris is the diaphragm that serves as the aperture stop. Refraction in the cornea causes
the effective aperture (the entrance pupil) to differ slightly from the physical pupil diameter. The entrance pupil is typically about 4
mm in diameter, although it can range from 2 mm (f/8.3) in a brightly lit place to 8 mm (f/2.1) in the dark. The latter value
decreases slowly with age; older people’s eyes sometimes dilate to not more than 5-6mm.
The approximate field of view of an individual human eye is 95° away from the nose, 75° downward, 60° toward the nose, and 60°
upward, allowing humans to have an almost 180-degree forward-facing horizontal field of view. With eyeball rotation of about 90°
(head rotation excluded, peripheral vision included), horizontal field of view is as high as 170°. About 12–15° temporal and 1.5°
below the horizontal is the optic nerve or blind spot which is roughly 7.5° high and 5.5° wide.
learning objectives
Identify factors responsible for nearsightedness and farsightedness vision defects
The human eye is the gateway to one of our five senses. The human eye is an organ that reacts with light. It allows light perception,
color vision, and depth perception, but not all eyes are perfect. A normal human eye can see about 10 million different colors!
Properties
Contrary to what you might think, the human eye is not a perfect sphere, but is made up of two differently shaped pieces, the
cornea and the sclera. These two parts are connected by a ring called the limbus. The part of the eye that is seen is the iris, which is
25.1.4 https://phys.libretexts.org/@go/page/16186
the colorful part of the eye. In the middle of the iris is the pupil, which is the black dot that changes size. The cornea covers these
elements, but it is transparent. The fundus is on the opposite of the pupil, but inside the eye and cannot be seen without special
instruments. The optic nerve is what conveys the signals of the eye to the brain. shows a diagram of the eye.
Vision
The different parts of the eye have different refractive indexes, and this is what bends the rays to form an image. The cornea
provides two-thirds of the power to the eye. The lens provides the remaining power. The image passes through several layers of the
eye, but this happens in a way very similar to that of a convex lens. When the image finally reaches the retena, it is inverted, but the
brain will correct this. For the vision to be clear, the image has to be formed directly on the retina. The focus needs to be changed,
much like a camera, depending on the distance and size of the object. The eye’s lens is flexible, and changes shape. This changes
the focal length. The eye’s ciliary muscles control the shape of the lens. When you focus on something, you squeeze or relax these
muscles.
Vision Diagram: An image is formed on the retina with light rays converging most at the cornea and upon entering and exiting the
lens. Rays from the top and bottom of the object are traced and produce an inverted real image on the retina. The distance to the
object is drawn smaller than scale.
Near Sighted Vision: This occurs when the image is formed before the retina
Far Sighted Vision: This occurs when the image is formed behind the retina
25.1.5 https://phys.libretexts.org/@go/page/16186
Key Points
The eye is made up of a number of parts, including the iris, pupil, cornea, and retina.
The eye has six muscles which control the eye movement, all providing different tension and torque.
The eye works a lot like a camera, the pupil provides the f-stop, the iris the aperture stop, the cornea resembles a lens. The way
that the image is formed is much like the way a convex lens forms an image.
The cones in the retina are responsible for seeing colors. There are three types of cones, each type can only pick up one color:
red, green, or blue. This is why TVs and computer screens are made up of thousands of little red, green, or blue lights.
The human eye is more sensitive to intensity changes than color changes, which is why it is acceptable to use black and white
photography in place of color and why people can still distinguish everything in the photo without colors.
Colors are usually written in different values of red, green, or blue. Each value is the log form of that frequency.
The retina of human eye has a static contrast ratio of around 100:1 and a dynamic contrast ratio of about 1,000,000:1.
The eye includes a lens not dissimilar to lenses found in optical instruments, such as cameras.
The approximate field of view of an individual human eye is 95° away from the nose, 75° downward, 60° toward the nose, and
60° upward, allowing humans to have an almost 180-degree forward-facing horizontal field of view.
The focal point of the image will change depending on how the lens is shaped. Your lens changes depending on the distance of
the object, by the relaxation or contraction of the muscles, and this controls the focal length.
Near sightedness occurs when the image is formed before the retina.
Farsightedness occurs when the image is formed behind the retina.
Key Terms
pupil: The hole in the middle of the iris of the eye, through which light passes to be focused on the retina.
aperture: The diameter of the aperture that restricts the width of the light path through the whole system. For a telescope, this is
the diameter of the objective lens (e.g., a telescope may have a 100 cm aperture).
luminance: The intensity of an object, independent from its color.
static contrast ratio: Luminosity ratio of the brightest and darkest color the system is capable of processing simultaneously at
any instant of time.
dynamic contrast ratio: Luminosity ratio of the brightest and darkest color the system is capable of processing over time
(while the picture is moving).
field of view: The angular extent of what can be seen, either with the eye or with an optical instrument or camera.
myopia: A disorder of the vision where distant objects appear blurred because the eye focuses their images in front of the retina
instead of on it.
hyperopia: A disorder of the vision where the eye focusses images behind the retina instead of on it, so that distant objects can
be seen better than near objects.
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25.2: Other Optical Instruments
A magnifying glass is a convex lens that lets the observer see a larger image of the object under observation. The lens is usually
mounted in a frame with a handle, as shown below.
Magnifying Glass: A magnifying glass is a convex lens that lets the observer see a larger image of the object under observation.
The magnification of a magnifying glass depends upon where the instrument is placed between the user’s eye and the object being
viewed and upon the total distance between eye and object. The magnifying power is the ratio of the sizes of the images formed on
the user’s retina with and without the magnifying glass. When not using the lens, the user would typically bring the object as close
to the eye as possible without it becoming blurry. (This point, known as the near point, varies with age. In a young child its distance
can be as short as five centimeters, while in an elderly person its distance may be as long as one or two meters. ) Magnifiers are
typically characterized using a “standard” value of 0.25m.
The highest magnifying power is obtained by putting the lens very close to the eye and moving both the eye and the lens together to
obtain the best focus. When the lens is used this way, the magnifying power can be found with the following equation:
1
MP 0 = ⋅ Φ +1 (25.2.1)
4
where Φ= optical power. When the magnifying glass is held close to the object and the eye is moved away, the magnifying power
is approximated by:
1
MP 0 = ⋅Φ (25.2.2)
4
Typical magnifying glasses have a focal length of 25cm and an optical power of four diopters. This type of glass would be sold as a
2x magnifier, but a typical observer would see about one to two times magnification depending on the lens position.
The earliest evidence of a magnifying device was Aristophanes’s “lens” from 424 BC, a glass globe filled with water. (Seneca
wrote that it could be used to read letters “no matter how small or dim.”) Roger Bacon described the properties of magnifying
glasses in the 13th century, and eyeglasses were also developed in 13th-century Italy.
The Camera
Cameras are optical devices that allow a user to record an image of an object, either on photo paper or digitally.
What is a Camera?
A camera is a device that allows you to record images, either on film or digitally. Cameras can record images as well as movies;
movies themselves got their name from moving pictures. The word camera comes from the Latin phrase camera obscura, which
means “dark chamber. ” The camera obscura was an early instrument for projecting images from slides. The camera that you use
today is an evolution of the camera obscura.
A camera is usually comprised of an opening, or aperture, that allows light to enter into a hollow area and a surface that records the
light at the other end. In the 20th century, these images would be stored on photographic paper that then had to be developed, but
now most cameras store images digitally.
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How does a Camera Work?
Cameras have a lot of components that allow them to work. Let’s look at them one at a time.
The Lens
The camera lens allows the light to enter into the camera and is typically convex. There are many types of lenses that can be used,
each for a different type of photography. There are lenses for close-ups, for sports, for architecture, and for portraits.
The two major features of a lens are focal length and aperture. The focal length determines the magnification of the image, and the
aperture controls the light intensity. The f-number on a camera controls the shutter speed. This is the speed at which the shutter,
which acts as its “eyelid,” opens and closes. The larger the aperture, the smaller the f-number must be in order to get the shutter
opened and closed fully. The time it takes to open and close the shutter is called the exposure. shows an example of two lenses of
the same size but with different apertures.
Focus
Some cameras have a fixed focus, and only objects of a certain size at a certain distance from the camera will be in focus. Other
cameras allow you to manually or automatically adjust the focus. shows a picture taken with a camera with manual focus; this
allows the user to determine which objects will be in focus and which will not. The range of distance within which objects appear
sharp and clear is called the depth of field.
Exposure
The aperture controls the intensity of the light entering the camera, and the shutter controls the exposure — the amount of time that
the light is allowed into the camera.
Shutter
The shutter is what opens and closes to allow light through the aperture. The speed at which it opens and closes is called the f-
number. For a larger aperture, the f-number is generally small for a quick shutter speed. For a smaller aperture, the f-number is
larger, allowing for a slower shutter speed.
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Diagram of a Compound Microscope: This diagram shows the setup of mirrors that allow for the magnification of images.
Since each lens produces a magnification that multiplies the height of the image, the total magnification is a product of the
individual magnifications. The equation for calculating this is as follows:
m = mo me (25.2.3)
The Telescope
The telescope aids in observation of remote objects by collecting electromagnetic radiation, such as visible light.
The telescope aids in observation of remote objects by collecting electromagnetic radiation, such as x-rays, visible light, infrared,
and submillimeter rays. The first telescopes were invented in the Netherlands in the 1600s and used glass lenses. Shortly after,
people began to build them using mirrors and called them reflecting telescopes.
History
The first telescope was a refracting telescope made by spectacle makers in the Netherlands in 1608. In 1610, Galileo made his own
improved design. After the refracting telescope was invented, people began to explore the idea of a telescope that used mirrors. The
potential advantages of using mirrors instead of lenses were a reduction in spherical aberrations and the elimination of chromatic
aberrations. In 1668, Newton built the first practical reflecting telescope. With the invention of achromatic lenses in 1733, color
aberrations were partially corrected, and shorter, more functional refracting telescopes could be constructed. Reflecting telescopes
were not practical because of the highly corrosive metals used to make mirrors until the introduction of silver-coated glass mirrors
in 1857.
Types of Telescopes
Refracting Telescopes
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Schematic of Keplerian Refracting Telescope: All refracting telescopes use the same principles. The combination of an objective
lens 1 and some type of eyepiece 2 is used to gather more light than the human eye is able to collect on its own, focus it 5, and
present the viewer with a brighter, clearer, and magnified virtual image 6.
The figure above is a diagram of a refracting telescope. The objective lens (at point 1) and the eyepiece (point 2) gather more light
than a human eye can collect by itself. The image is focused at point 5, and the observer is shown a brighter, magnified virtual
image at point 6. The objective lens refracts, or bends, light. This causes the parallel rays to converge at a focal point, and those
that are not parallel converge on a focal plane.
Reflecting Telescopes
Reflecting telescopes, such as the one shown in, use either one or a combination of curved mirrors that reflect light to form an
image. They allow an observer to view objects that have very large diameters and are the primary type of telescope used in
astronomy. The object being observed is reflected by a curved primary mirror onto the focal plane. (The distance from the mirror to
the focal plane is called the focal length. ) A sensor could be located here to record the image, or a secondary mirror could be added
to redirect the light to an eyepiece.
Catadioptric Telescopes
Catadioptric telescopes, such as the one shown in, combine mirrors and lenses to form an image. This system has a greater degree
of error correction than other types of telescopes. The combination of reflective and refractive elements allows for each element to
correct the errors made by the other.
X-Ray Diffraction
The principle of diffraction is applied to record interference on a subatomic level in the study of x-ray crystallography.
X-ray diffraction was discovered by Max von Laue, who won the Nobel Prize in physics in 1914 for his mathematical evaluation of
observed x-ray diffraction patterns.
Diffraction is the irregularities caused when waves encounter an object. You have most likely observed the effects of diffraction
when looking at the bottom of a CD or DVD. The rainbow pattern that appears is a result of the light being interfered by the pits
and lands on the disc that hold the data. shows this effect. Diffraction can happen to any type of wave, not just visible light waves.
Bragg Diffraction
In x-ray crystallography, the term for diffraction is Bragg diffraction, which is the scattering of waves from a crystalline structure.
William Lawrence Bragg formulated the equation for Bragg’s law, which relates wavelength to the angle of incidence and lattice
spacing. Refer to for a diagram of the following equation: nλ=2dsin(θ)nλ=2dsin(θ)
n – numeric constant known as the order of the diffracted beam
λ – wavelength
d – distance between lattice planes
θ – angle of diffracted wave
The waves will experience either constructive interference or destructive interference. Similarly, the x-ray beam that is diffracted
off a crystal will have some parts that have stronger energy, and others that lose energy. This depends on the wavelength and the
lattice spacing.
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The X-ray Diffractometer
The XRD machine uses copper metal as the element for the x-ray source. Diffraction patterns are recorded over an extended period
of time, so it is very important that the beam intensity remains constant. Film used to be used to record the data, but that was
inconvenient because it had to be replaced often. Now the XRD machines are equipped with semiconductor detectors. These XRD
machines record images in two ways, either continuous scans or step scanning. In continuous scans, the detector moves in circular
motions around the object, while a beam of x-ray is constantly shot at the detector. Pulses of energy are plotted with respect to
diffraction angle. The step scan method is the more popular method. It is much more efficient than continuous scans. In this
method, the detector collects data at a single fixed angle at a time. To ensure that the incident beam is continuous, XRD machines
are equipped with a Soller slit. This acts like polarized sunglasses by organizing random x-ray beams into a stack of neatly
arranged waves parallel to the plane of the detector.
Overview
X-ray imaging, or radiography, used x-rays to view material within the body that cannot be seen by the human eye by identifying
areas of different density and composition. CT Scans use the assistance of a computer to take this information, and generate 3
dimensional images.
X-ray Imaging
X-ray radiographs are produced by projecting a beam of X-rays toward an object, in medical cases, a part of the human body.
Depending on the physical properties of the object (density and composition), some of the X-rays can be partially absorbed. The
portion of the rays that are not absorbed then pass through the object and are recorded by either film or a detector, like in a camera.
This provides the observer with a 2 dimensional representation of all the components of that object superimposed on each other.
shows an image of a human elbow.
Tomography
Tomography refers to imaging by sections, or sectioning. demonstrates this concept. The three-dimensional image is broken down
into sections. (S1) shows a section from the left and (S2) shows a section from the right.
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CT Scans
CT scans, or computed tomography scans use a combination of X-ray radiography and tomography to produce slices of areas of the
human body. Doctors can analyze the area, and based on the ability of the material to block the X-ray beam, understand more about
the material. shows a CT Scan of a human brain. Doctors can cross reference the images with known properties of the same
material and determine if there are any inconsistencies or problems. Although generally these scans are shown as in, the
information recorded can be used to create a 3 dimensional image of the area. shows a three dimensional image of a brain that was
made by compiling CT Scans.
Microscope Classes
One way to group microscopes is based on how the image is generated through the microscope. Here are three ways we can
classify microscopes:
1. ) Light or Photon – optical microscopes
2. ) Electrons – electron microscopes
3. ) Probe – scanning probe microscopes.
Microscopes can also be classified based on whether they analyze the sample by scanning a point at a time (scanning electron
microscopes), or by analyzing the entire sample at once (transmission electron microscopes).
Types of Microscopes
In optical microscopes, the better the contrast between the image and the surface it is being viewed on, the better the resolution
will be to the viewer. There are many illumination techniques to generate improved contrast. These techniques include “dark
field” and “bright field.” With the dark field technique the light is scattered by the object and the image appears to the observer
on a dark background. With the bright field technique the object is illuminated from below to increase the contrast in the image
seen by viewers.
Transmission Electron Microscope: The TEM passes electrons through the sample, and allows people to see objects that are
normally not seen by the naked eye. A beam of electrons is transmitted through an ultra thin specimen, interacting with the
specimen as it passes through. This interaction forms an image that is magnified and focused onto an imaging device.
Scanning Electron Microscopes: Referred to as SEM, these microscopes look at the surface of objects by scanning them with a
fine electron beam. The electron beam of the microscope interacts with the electrons in the sample and produces signals that can
be detected and have information about the topography and composition.
Atomic Force Microscopy: The AFM is a scanning probe type of microscopy with very high resolution and is one of the
foremost tools for imaging at the nanoscale. The mechanical probe feels the surface with a cantilever with a sharp tip. The
deflection of the tip is then measured using a laser spot that is reflected from the surface of the cantilever.
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object. Since effects of diffraction become most prominent for waves whose wavelength is roughly similar to the dimensions of the
diffracting objects, the wavelength of the imaging beam sets a fundamental limit on the resolution of any optical system.
Airy Disk: Computer-generated image of an Airy disk. The gray scale intensities have been adjusted to enhance the brightness of
the outer rings of the Airy pattern.
The denominator nsinθnsinθ is called the numerical aperture and can reach about 1.4 in modern optics, hence the Abbe limit is
roughly d=λ/2. With green light around 500 nm, the Abbe limit is 250 nm which is large compared to most nanostructures, or
biological cells with sizes on the order of 1μm and internal organelles which are much smaller. Using a 500 nm beam, you cannot
(in principle) resolve any features with size less than around 250 nm.
Improving Resolution
To increase the resolution, shorter wavelengths can be used such as UV and X-ray microscopes. These techniques offer better
resolution but are expensive, suffer from lack of contrast in biological samples and may damage the sample. There are techniques
for producing images that appear to have higher resolution than allowed by simple use of diffraction-limited optics. Although these
techniques improve some aspect of resolution, they generally involve an enormous increase in cost and complexity. Usually the
technique is only appropriate for a small subset of imaging problems.
Aberrations
An aberration, or distortion, is a failure of rays to converge at one focus because of limitations or defects in a lens or mirror.
Chromatic Aberration
A chromatic aberration, also called achromatism or chromatic distortion, is a distortion of colors. This aberration happens when the
lens fails to focus all the colors on the same convergence point. This happens because lenses have a different index of refraction for
different wavelengths of light. The refractive index decreases with increasing wavelength. These aberrations or distortions occur on
the edges of color boundaries between bright and dark areas of an image. Since the index of refraction of lenses depends on color
or wavelength, images are produced at different places and with different magnifications for different colors. shows chromatic
aberration for a single convex lens. Since violet rays have a higher refractive index than red, they are bent more and focused closed
to the lens. shows a two-lens system using a diverging lens to partially correct for this, but it is nearly impossible to do so
completely.
The law of reflection is independent of wavelength, and therefore mirrors do not have this problem. This is why it is advantageous
to use mirrors in telescopes and other optical systems.
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Comatic Aberration
A comatic aberration, or coma, occurs when the object is off-center. Different parts of a lens of a mirror do not refract or reflect the
image to the same point, as shown in. They can also be result of an imperfection in the lens or other component and result in off-
axis point sources. These aberrations can cause objects to appear pear-shaped. They can also cause stars to appear distorted or
appear to have tails, as with comets.
Other Aberrations
Spherical aberrations are a form of aberration where rays converging from the outer edges of a lens converge to a focus closer to
the lens, and rays closer to the axis focus further. Astigmatisms are also a form of aberration in the lenses of the eyes where rays
that propagate in two perpendicular planes have different foci. This can eventually cause a monochromatic image to distort
vertically or horizontally. Another aberration or distortion is a barrel distortion where image magnification decreases with the
distance from the optical axis. The apparent effect is that of an image which has been mapped around a sphere, like in a fisheye
lens.
Key Points
The magnification of a magnifying glass depends upon where it is placed between the user’s eye and the object being viewed
and upon the total distance between eye and object.
The magnifying power is the ratio of the sizes of the images formed on the user’s retina with and without the lens.
The highest magnifying power is obtained by putting the lens very close to the eye and moving both the eye and the lens
together to obtain the best focus.
Cameras work very similarly to how the human eye works. The iris is similar to the lens; the pupil is similar to the aperture; and
the eyelid is similar to the shutter.
Cameras are a modern evolution of the camera obscura. The camera obscura was a device used to project images.
The most important part of a camera is the lens, which allows the image to be magnified and focused. This can be done
manually on some cameras and automatically on newer cameras.
Movie cameras work by taking many pictures each second and then showing each image in order very quickly to give the effect
that the pictures are moving. This is where the name “movie” comes from.
A compound microscope uses multiple lenses to create an enlarged image that is easier for a human eye to see; this is due to the
fact that the final image is farther away from the observer, and therefore the eye is more relaxed when viewing the image.
The object is placed just beyond to focal length of the objective lens. An enlarged image is then captured by the objective lens,
which acts as the object for the ocular lens. The ocular lens is closer to the new image than its focal length, causing it to act as a
magnifying glass.
Since the final image is just a multiple of the size of the first image, the final magnification is a product of both magnifications
from each lens.
Until the invention of silver-backed mirrors, refractive mirrors were the standard for use in telescopes. This was because of the
highly corrosive nature of the metals used in older mirrors. Since then, reflective mirrors have replaced refractive mirrors in
astronomy.
There are three main types of optical telescopes: refractive, reflective, and catadioptric.
Refractive telescopes, such as the one invented by Galileo, use an objective lens and an eyepiece. The image is focused at the
focal point and allows the observer to see a brighter, larger image than they would with their own eye.
Reflective telescopes use curved mirrors that reflect light to form an image. Sometimes a secondary mirror redirects the image
to an eyepiece. Other times the image is recorded by a sensor and observed on a computer screen.
Catadioptric telescopes combine mirrors and lenses to form an image. This system has a greater degree of error correction than
other types of telescopes. The combination of reflective and refractive elements allows for each element to correct the errors
made by the other.
Diffraction is what happens when waves encounter irregularities on a surface or object and are caused to interfere with each
other, either constructively or destructively.
The Bragg law pertains to applying the laws of diffraction to crystallography in order to obtain precise images of the lattice
structures in atoms.
The x-ray diffractometer is the machine used to scan the object by shooting a wave at it and recording the interference it
encounters.
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Most XRDs are equipped with a Soller slit, which acts like a polarizer for the incident beam. It makes sure that the incident
beam being recorded is perfectly parallel to the object being analyzed.
Radiography uses x-rays to take pictures of materials with in an object that can not be seen. They shoot x-ray beams through the
object, and collect the rays on film or a detector on the other side. Some of the rays are absorbed into the denser materials, and
this is how the image is produced.
X-ray radiographs take images of all the materials within an object superimposed on each other.
The traditional, superimposed images can be helpful for a number of applications, but CT scans enable the observer to see just
the desired sections of a material.
Modern CT scans can even take all the slices, or layers, and arrange them into a three-dimensional representation of the object.
For better resolution, it is important that there is a lot of contrast between the image and the background.
Microscopes are classified by what interacts with the object, such as light or electrons. They are also classified by whether they
take images by scanning a little at a time or by taking images of the entire object at once.
Some common types of specialty microscopes are scanning electron micrscopes (SEM), transmission electron microscopes
(TEM), both of which are electron microscopes, and atomic force microscopes (ATM) which is a scanning probe microscope.
Since effects of diffraction become most prominent for waves whose wavelength is roughly similar to the dimensions of the
diffracting objects, the wavelength of the imaging beam sets a fundamental limit on the resolution of any optical system.
The Abbe diffraction limit for a microscope is given as d = λ
2(n sin θ)
.
Since the diffraction limit is proportional to wavelength, to increase the resolution, shorter wavelengths can be used such as UV
and X-ray microscopes.
There are many types of aberrations, including chromatic, spherical, comatic, astigmatism, and barrel distortion.
Chromatic aberrations occur due to the fact that lenses have different refractive indexes for different wavelengths, and therefore
colors. These aberrations occur right on the edges of images between light and dark areas of the picture.
Mirrors do not have chromatic aberrations because they do not rely on the index of refraction, but rather the index of reflection,
which is independent of wavelength.
Comatic aberrations are due to imperfections in lenses and cause the point source to be off-center. This can cause images to
appear pear-shaped, or cause images to have tails, as with comets.
Key Terms
lens: an object, usually made of glass, that focuses or defocuses the light that passes through it
diopter: a unit of measure of the power of a lens or mirror, equal to the reciprocal of its focal length in meters. Myopia is
diagnosed and measured in diopters
convex: curved or bowed outward like the outside of a bowl or sphere or circle
shutter speed: The duration of time for which the shutter of a camera remains open when exposing photographic film or other
photosensitive material to light for the purpose of recording an image
chromatic aberration: an optical aberration, in which an image has colored fringes, caused by differential refraction of light of
different wavelengths
spherical aberration: a type of lens aberration that causes blurriness, particularly away from the center of the lens
achromatic: free from color; transmitting light without color-related distortion
constructive interference: Occurs when waves interfere with each other crest to crest and the waves are exactly in phase with
each other.
crystallography: The experimental science of determining the arrangement of atoms in solids.
destructive interference: Occurs when waves interfere with each other crest to trough (peak to valley) and are exactly out of
phase with each other.
radiography: The use of X-rays to view a non-uniformly composed material such as the human body.
tomography: Imaging by sections or sectioning.
superimposed: Positioned on or above something else, especially in layers
contrast: A difference in lightness, brightness and/or hue between two colors that makes them more or less distinguishable
diffraction: The bending of a wave around the edges of an opening or an obstacle.
nanostructure: Any manufactured structure having a scale between molecular and microscopic.
aperture: The diameter of the aperture that restricts the width of the light path through the whole system. For a telescope, this is
the diameter of the objective lens (e.g., a telescope may have a 100 cm aperture).
distortion: (optics) an aberration that causes magnification to change over the field of view.
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refraction: Changing of a light ray’s direction when it passes through variations in matter.
aberration: The convergence to different foci, by a lens or mirror, of rays of light emanating from one and the same point, or
the deviation of such rays from a single focus; a defect in a focusing mechanism that prevents the intended focal point.
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CHAPTER OVERVIEW
26: Wave Optics is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
26.1: Superposition and Interference
learning objectives
Contrast the effects of constructive and destructive interference
Wave Interference: Examples of constructive (left) and destructive (right) wave interference.
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Wave Interference: A brief introduction to constructive and destructive wave interference and the principle of superposition.
A simple form of wave interference is observed when two waves of the same frequency (also called a plane wave) intersect at an
angle, as shown in. Assuming the two waves are in phase at point B, then the relative phase changes along the x-axis. The phase
difference at point A is given by:
Interference of Plane Waves: Geometrical arrangement for two plane wave interference.
2πd 2πx sin θ
Δφ = = (26.1.1)
λ λ
Destructive interference occurs when the waves are half a cycle out of phase, or
x sin θ 1 3
=± ,± ,… (26.1.3)
λ 2 2
Air Wedge
An air wedge is a simple interferometer used to visualize the disturbance of the wave front after propagation through a test object.
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learning objectives
Describe how an air wedge is used to visualize the disturbance of a wave front after proagation
Air Wedge
An air wedge is one of the simplest designs of shearing interferometers used to visualize the disturbance of the wave front after
propagation through a test object. An air wedge can be used with nearly any light source, including non-coherent white light. The
interferometer consists of two optical glass wedges (~2-5 degrees), pushed together and then slightly separated from one side to
create a thin air-gap wedge. An example of an air wedge interferometer is shown in.
Light Reflections Inside an Air Wedge Interferometer: Beam path inside of air wedge interferometer
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Newton’s Rings
Newton’s rings are a series of concentric circles centered at the point of contact between a spherical and a flat surface.
learning objectives
Apply Newton’s rings to determine light characteristics of a lens
Newton’s Rings
In 1717, Isaac Newton first analyzed an interference pattern caused by the reflection of light between a spherical surface and an
adjacent flat surface. Although first observed by Robert Hooke in 1664, this pattern is called Newton’s rings, as Newton was the
first to analyze and explain the phenomena. Newton’s rings appear as a series of concentric circles centered at the point of contact
between the spherical and flat surfaces. When viewed with monochromatic light, Newton’s rings appear as alternating bright and
dark rings; when viewed with white light, a concentric ring pattern of rainbow colors is observed. An example of Newton’s rings
when viewed with white light is shown in the figure below.
Newton’s Rings in a drop of water: Newton’s rings seen in two plano-convex lenses with their flat surfaces in contact. One
surface is slightly convex, creating the rings. In white light, the rings are rainbow-colored, because the different wavelengths of
each color interfere at different locations.
The light rings are caused by constructive interference between the light rays reflected from both surfaces, while the dark rings are
caused by destructive interference. The outer rings are spaced more closely than the inner ones because the slope of the curved lens
surface increases outwards. The radius of the Nth bright ring is given by:
1/2
1
rN = [(N − λR)] (26.1.4)
2
where N is the bright-ring number, R is the radius of curvature of the lens the light is passing through, and λ is the wavelength of
the light passing through the glass.
A spherical lens is placed on top of a flat glass surface. An incident ray of light passes through the curved lens until it comes to the
glass-air boundary, at which point it passes from a region of higher refractive index n (the glass) to a region of lower n (air). At this
boundary, some light is transmitted into the air, while some light is reflected. The light that is transmitted into the air does not
experience a change in phase and travels a a distance, d, before it is reflected at the flat glass surface below. This second air-glass
boundary imparts a half-cycle phase shift to the reflected light ray because air has a lower n than the glass. The two reflected light
rays now travel in the same direction to be detected. As one gets farther from the point at which the two surfaces touch, the distance
dincreases because the lens is curving away from the flat surface.
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Formation of Interference Fringes: This figure shows how interference fringes form.
If the path length difference between the two reflected light beams is an odd multiple of the wavelength divided by two, λ/2, the
reflected waves will be 180 degrees out of phase and destructively interfere, causing a dark fringe. If the path-length difference is
an even multiple of λ/2, the reflected waves will be in phase with one another. The constructive interference of the two reflected
waves creates a bright fringe.
Key Points
When two or more waves are incident on the same point, the total displacement at that point is equal to the vector sum of the
displacements of the individual waves.
Light exhibits wave characteristics in various media as well as in a vacuum. When light goes from a vacuum to some medium,
like water, its speed and wavelength change, but its frequency f remains the same.
When light is reflected off a medium with a higher index of refraction, crests get reflected as troughs and troughs get reflected
as crests. In other words, the wave undergoes a 180 degree change of phase upon reflection, and the reflected ray “jumps” ahead
by half a wavelength.
An air wedge interferometer consists of two optical glass wedges (~2-5 degrees), pushed together and then slightly separated
from one side to create a thin air-gap wedge.
The air gap between the two glass plates has two unique properties: it is very thin (micrometer scale) and has perfect flatness.
To minimize image aberrations of the resulting fringes, the angle plane of the glass wedges has to be placed orthogonal to the
angle plane of the air wedge.
When viewed with monochromatic light, Newton’s rings appear as alternating bright and dark rings; when viewed with white
light, a concentric ring pattern of rainbow colors is observed.
If the path length difference between the two reflected light beams is an odd multiple of the wavelength divided by two, λ/2, the
reflected waves will be 180 degrees out of phase and destructively interfere, causing a dark fringe.
If the path length difference is an even multiple of λ/2, the reflected waves will be in-phase with one another. The constructive
interference of the two reflected waves creates a bright fringe.
Key Terms
coherent: Of waves having the same direction, wavelength and phase, as light in a laser.
plane wave: A constant-frequency wave whose wavefronts (surfaces of constant phase) are infinite parallel planes of constant
peak-to-peak amplitude normal to the phase velocity vector.
orthogonal: Of two objects, at right angles; perpendicular to each other.
interferometer: Any of several instruments that use the interference of waves to determine wavelengths and wave velocities,
determine refractive indices, and measure small distances, temperature changes, stresses, and many other useful measurements.
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ, and corresponds to the velocity of the wave divided by its frequency.
lens: an object, usually made of glass, that focuses or defocuses the light that passes through it
monochromatic: Describes a beam of light with a single wavelength (i.e., of one specific color or frequency).
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26.2: Diffraction
learning objectives
Formulate Huygens’s Principle
Overview
The Huygens-Fresnel principle states that every point on a wavefront is a source of wavelets. These wavelets spread out in the
forward direction, at the same speed as the source wave. The new wavefront is a line tangent to all of the wavelets.
Background
Christiaan Huygens was a Dutch scientist who developed a useful technique for determining how and where waves propagate. In
1678, he proposed that every point that a luminous disturbance touches becomes itself a source of a spherical wave. The sum of the
secondary waves (waves that are a result of the disturbance) determines the form of the new wave. shows secondary waves
traveling forward from their point of origin. He was able to come up with an explanation of the linear and spherical wave
propagation, and derive the laws of reflection and refraction (covered in previous atoms ) using this principle. He could not,
however, explain what is commonly known as diffraction effects. Diffraction effects are the deviations from rectilinear propagation
that occurs when light encounters edges, screens and apertures. These effects were explained in 1816 by French physicist Augustin-
Jean Fresnel.
Straight Wavefront: Huygens’s principle applied to a straight wavefront. Each point on the wavefront emits a semicircular
wavelet that moves a distance s=vt. The new wavefront is a line tangent to the wavelets.
Huygens’s Principle
Figure 1 shows a simple example of the Huygens’s Principle of diffraction. The principle can be shown with the equation below:
s = vt (26.2.1)
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Huygens’s Refraction: Huygens’s principle applied to a straight wavefront traveling from one medium to another where its speed
is less. The ray bends toward the perpendicular, since the wavelets have a lower speed in the second medium.
Reflection: Huygens’s principle applied to a straight wavefront striking a mirror. The wavelets shown were emitted as each point
on the wavefront struck the mirror. The tangent to these wavelets shows that the new wavefront has been reflected at an angle equal
to the incident angle. The direction of propagation is perpendicular to the wavefront, as shown by the downward-pointing arrows.
Example 26.2.1:
This principle is actually something you have seen or experienced often, but just don’t realize. Although this principle applies
to all types of waves, it is easier to explain using sound waves, since sound waves have longer wavelengths. If someone is
playing music in their room, with the door closed, you might not be able to hear it while walking past the room. However, if
that person where to open their door while playing music, you could hear it not only when directly in front of the door opening,
but also on a considerable distance down the hall to either side. is a direct effect of diffraction. When light passes through
much smaller openings, called slits, Huygens’s principle shows that light bends similar to the way sound does, just on a much
smaller scale. We will examine in later atoms single slit diffraction and double slit diffraction, but for now it is just important
that we understand the basic concept of diffraction.
Diffraction
As we explained in the previous paragraph, diffraction is defined as the bending of a wave around the edges of an opening or an
obstacle.
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learning objectives
Explain why Young’s experiment more credible than Huygens’
The double-slit experiment, also called Young’s experiment, shows that matter and energy can display both wave and particle
characteristics. As we discussed in the atom about the Huygens principle, Christiaan Huygens proved in 1628 that light was a
wave. But some people disagreed with him, most notably Isaac Newton. Newton felt that color, interference, and diffraction effects
needed a better explanation. People did not accept the theory that light was a wave until 1801, when English physicist Thomas
Young performed his double-slit experiment. In his experiment, he sent light through two closely spaced vertical slits and observed
the resulting pattern on the wall behind them. The pattern that resulted can be seen in.
Young’s Double Slit Experiment: Light is sent through two vertical slits and is diffracted into a pattern of vertical lines spread out
horizontally. Without diffraction and interference, the light would simply make two lines on the screen.
Wave-Particle Duality
The wave characteristics of light cause the light to pass through the slits and interfere with itself, producing the light and dark areas
on the wall behind the slits. The light that appears on the wall behind the slits is scattered and absorbed by the wall, which is a
characteristic of a particle.
Young’s Experiment
Why was Young’s experiment so much more credible than Huygens’? Because, while Huygens’ was correct, he could not
demonstrate that light acted as a wave, while the double-slit experiment shows this very clearly. Since light has relatively short
wavelengths, to show wave effects it must interact with something small — Young’s small, closely spaced slits worked.
The example in uses two coherent light sources of a single monochromatic wavelength for simplicity. (This means that the light
sources were in the same phase. ) The two slits cause the two coherent light sources to interfere with each other either
constructively or destructively.
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Practical Constructive and Destructive Wave Interference: Double slits produce two coherent sources of waves that interfere.
(a) Light spreads out (diffracts) from each slit because the slits are narrow. These waves overlap and interfere constructively (bright
lines) and destructively (dark regions). We can only see this if the light falls onto a screen and is scattered into our eyes. (b)
Double-slit interference pattern for water waves are nearly identical to that for light. Wave action is greatest in regions of
constructive interference and least in regions of destructive interference. (c) When light that has passed through double slits falls on
a screen, we see a pattern such as this.
Theoretical Constructive and Destructive Wave Interference: The amplitudes of waves add together. (a) Pure constructive
interference is obtained when identical waves are in phase. (b) Pure destructive interference occurs when identical waves are
exactly out of phase (shifted by half a wavelength).
The pattern that results from double-slit diffraction is not random, although it may seem that way. Each slit is a different distance
from a given point on the wall behind it. For each different distance, a different number of wavelengths fit into that path. The
26.2.4 https://phys.libretexts.org/@go/page/16195
waves all start out in phase (matching crest-to-crest), but depending on the distance of the point on the wall from the slit, they could
be in phase at that point and interfere constructively, or they could end up out of phase and interfere with each other destructively.
learning objectives
Describe function of the diffraction grating
Diffraction Grating
A diffraction grating is an optical component with a periodic structure that splits and diffracts light into several beams travelling in
different directions. The directions of these beams depend on the spacing of the grating and the wavelength of the light so that the
grating acts as the dispersive element. Because of this, gratings are often used in monochromators, spectrometers, wavelength
division multiplexing devices, optical pulse compressing devices, and many other optical instruments.
A photographic slide with a fine pattern of purple lines forms a complex grating. For practical applications, gratings generally have
ridges or rulings on their surface rather than dark lines. Such gratings can be either transmissive or reflective. Gratings which
modulate the phase rather than the amplitude of the incident light are also produced, frequently using holography.
Ordinary pressed CD and DVD media are every-day examples of diffraction gratings and can be used to demonstrate the effect by
reflecting sunlight off them onto a white wall. (see ). This is a side effect of their manufacture, as one surface of a CD has many
small pits in the plastic, arranged in a spiral; that surface has a thin layer of metal applied to make the pits more visible. The
structure of a DVD is optically similar, although it may have more than one pitted surface, and all pitted surfaces are inside the
disc. In a standard pressed vinyl record when viewed from a low angle perpendicular to the grooves, one can see a similar, but less
defined effect to that in a CD/DVD. This is due to viewing angle (less than the critical angle of reflection of the black vinyl) and
the path of the light being reflected due to being changed by the grooves, leaving a rainbow relief pattern behind.
Readable Surface of a CD: The readable surface of a Compact Disc includes a spiral track wound tightly enough to cause light to
diffract into a full visible spectrum.
Some bird feathers use natural diffraction grating which produce constructive interference, giving the feathers an iridescent effect.
Iridescence is the effect where surfaces seem to change color when the angle of illumination is changed. An opal is another
example of diffraction grating that reflects the light into different colors.
X-Ray Diffraction
X-ray crystallography is a method of determining the atomic and molecular structure of a crystal, in which the crystalline atoms
cause a beam of X-rays to diffract into many specific directions. By measuring the angles and intensities of these diffracted beams,
a crystallographer can produce a three-dimensional picture of the density of electrons within the crystal. From this electron density,
the mean positions of the atoms in the crystal can be determined, as well as their chemical bonds, their disorder and various other
information.
In an X-ray diffraction measurement, a crystal is mounted on a goniometer and gradually rotated while being bombarded with X-
rays, producing a diffraction pattern of regularly spaced spots known as reflections (see ). The two-dimensional images taken at
different rotations are converted into a three-dimensional model of the density of electrons within the crystal using the
mathematical method of Fourier transforms, combined with chemical data known for the sample.
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Reflections in Diffraction Patterns: Each dot, called a reflection, in this diffraction pattern forms from the constructive
interference of scattered X-rays passing through a crystal. The data can be used to determine the crystalline structure.
learning objectives
Formulate the Huygens’s Principle
Diffraction
As we explained in a previous atom, diffraction is defined as the bending of a wave around the edges of an opening or obstacle.
Diffraction is a phenomenon all wave types can experience. It is explained by the Huygens-Fresnel Principle, and the principal of
superposition of waves. The former states that every point on a wavefront is a source of wavelets. These wavelets spread out in the
forward direction, at the same speed as the source wave. The new wavefront is a line tangent to all of the wavelets. The
superposition principle states that at any point, the net result of multiple stimuli is the sum of all stimuli.
Single Slit Diffraction – One Wavelength: Visualization of single slit diffraction when the slit is equal to one wavelength.
A slit that is wider than a single wave will produce interference -like effects downstream from the slit. It is easier to understand by
thinking of the slit not as a long slit, but as a number of point sources spaced evenly across the width of the slit. This can be seen in
Figure 2.
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Single Slit Diffraction – Four Wavelengths: This figure shows single slit diffraction, but the slit is the length of 4 wavelengths.
To examine this effect better, lets consider a single monochromatic wavelength. This will produce a wavefront that is all in the
same phase. Downstream from the slit, the light at any given point is made up of contributions from each of these point sources.
The resulting phase differences are caused by the different in path lengths that the contributing portions of the rays traveled from
the slit.
The variation in wave intensity can be mathematically modeled. From the center of the slit, the diffracting waves propagate
radially. The angle of the minimum intensity (θmin) can be related to wavelength (λ) and the slit’s width (d) such that:
The intensity (I) of waves at any angle can also be calculated as a relation to slit width, wavelength and intensity of the original
waves before passing through the slit:
2
sin(πx)
I(θ) = I0 ( ) (26.2.3)
πx
learning objectives
Explain meaning of the Rayleigh criterion
Resolution Limits
Along with the diffraction effects that we have discussed in previous atoms, diffraction also limits the detail that we can obtain in
images. shows three different circumstances of resolution limits due to diffraction:
Resolution Limits: (a) Monochromatic light passed through a small circular aperture produces this diffraction pattern. (b) Two
point light sources that are close to one another produce overlapping images because of diffraction. (c) If they are closer together,
they cannot be resolved or distinguished.
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(a) shows a light passing through a small circular aperture. You do not see a sharp circular outline, but a spot with fuzzy edges.
This is due to diffraction similar to that through a single slit.
(b) shows two point sources close together, producing overlapping images. Due to the diffraction, you can just barely
distinguish between the two point sources.
(c) shows two point sources which are so close together that you can no longer distinguish between them.
This effect can be seen with light passing through small apertures or larger apertures. This same effect happens when light passes
through our pupils, and this is why the human eye has limited acuity.
Rayleigh Criterion
In the 19th century, Lord Rayleigh invented a criteria for determining when two light sources were distinguishable from each other,
or resolved. According to the criteria, two point sources are considered just resolved (just distinguishable enough from each other
to recognize two sources) if the center of the diffraction pattern of one is directly overlapped by the first minimum of the diffraction
pattern of the other. If the distance is greater between these points, the sources are well resolved (i.e., they are easy to distingiush
from each other). If the distance is smaller, they are not resolved (i.e., they cannot be distinguished from each other). The equation
to determine this is:
λ
θ = 1.22 (26.2.5)
D
θ – angle the objects are separated by, in radian λ – wavelength of light D – aperture diameter. shows this concept visually. This
equation also gives the angular spreading of a source of light having a diameter D.
Rayleigh Criterion: (a) This is a graph of intensity of the diffraction pattern for a circular aperture. Note that, similar to a single
slit, the central maximum is wider and brighter than those to the sides. (b) Two point objects produce overlapping diffraction
patterns. Shown here is the Rayleigh criterion for being just resolvable. The central maximum of one pattern lies on the first
minimum of the other.
Key Points
Diffraction is the concept that is explained using Huygens’s Principle, and is defined as the bending of a wave around the edges
of an opening or an obstacle.
This principle can be used to define reflection, as shown in the figure. It can also be used to explain refraction and interference.
Anything that experiences this phenomenon is a wave. By applying this theory to light passing through a slit, we can prove it is
a wave.
The principle can be shown with the equation below: s=vt s – distance v – propagation speed t – time Each point on the
wavefront emits a wave at speed, v. The emitted waves are semicircular, and occur at t, time later. The new wavefront is tangent
to the wavelets.
The wave characteristics of light cause the light to pass through the slits and interfere with each other, producing the light and
dark areas on the wall behind the slits. The light that appears on the wall behind the slits is partially absorbed by the wall, a
characteristic of a particle.
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Constructive interference occurs when waves interfere with each other crest-to-crest and the waves are exactly in phase with
each other. Destructive interference occurs when waves interfere with each other crest-to-trough (peak-to-valley) and are
exactly out of phase with each other.
Each point on the wall has a different distance to each slit; a different number of wavelengths fit in those two paths. If the two
path lengths differ by a half a wavelength, the waves will interfere destructively. If the path length differs by a whole
wavelength the waves interfere constructively.
The directions of the diffracted beams depend on the spacing of the grating and the wavelength of the light so that the grating
acts as the dispersive element.
Gratings are commonly used in monochromators, spectrometers, wavelength division multiplexing devices, optical pulse
compressing devices, and other optical instruments.
Diffraction of X-ray is used in crystallography to produce the three-dimensional picture of the density of electrons within the
crystal.
The Huygens’s Principle states that every point on a wavefront is a source of wavelets. These wavelets spread out in the
forward direction, at the same speed as the source wave. The new wavefront is a line tangent to all of the wavelets.
If a slit is longer than a single wavelength, think of it instead as a number of point sources spaced evenly across the width of the
slit.
Downstream from a slit that is longer than a single wavelength, the light at any given point is made up of contributions from
each of these point sources. The resulting phase differences are caused by the different in path lengths that the contributing
portions of the rays traveled from the slit.
Diffraction plays a large part in the resolution at which we are able to see things. There is a point where two light sources can be
so close to each other that we cannot distinguish them apart.
When two light sources are close to each other, they can be: unresolved (i.e., not able to distinguish one from the other), just
resolved (i.e., only able to distinguish them apart from each other), and a little well resolved (i.e., easy to tell apart from one
another).
In order for two light sources to be just resolved, the center of one diffraction pattern must directly overlap with the first
minimum of the other diffraction pattern.
Key Terms
diffraction: The bending of a wave around the edges of an opening or an obstacle.
constructive interference: Occurs when waves interfere with each other crest to crest and the waves are exactly in phase with
each other.
destructive interference: Occurs when waves interfere with each other crest to trough (peak to valley) and are exactly out of
phase with each other.
interference: An effect caused by the superposition of two systems of waves, such as a distortion on a broadcast signal due to
atmospheric or other effects.
iridescence: The condition or state of being iridescent; exhibition of colors like those of the rainbow; a prismatic play of color.
diffraction: The bending of a wave around the edges of an opening or an obstacle.
monochromatic: Describes a beam of light with a single wavelength (i.e., of one specific color or frequency).
resolution: The degree of fineness with which an image can be recorded or produced, often expressed as the number of pixels
per unit of length (typically an inch).
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26.3: Further Topics
learning objectives
Describe process that leads to the thin film interference
Thin Film Interference: In this video I continue with my tutorials on Electromagnetism to Optics which is pitched at university
undergraduate level. I have intended for a long time to record videos which describe the transition made from classical
electromagnetism to optics. In many respects these videos will cover ‘wave’ optics. I devote much time to discussing the complex
26.3.1 https://phys.libretexts.org/@go/page/16196
exponential representation of waves, Maxwell’s Equations, the wave equation etc.Specifically here, I derive the formula for the
optical path difference and the phase difference for a ‘wave’ of light propagating through a thin film. This expression can be used
for anti reflective coatings. The phase difference is the product of the optical path differene and the wave vector k.I hope it’s of
use!! Thank you for watching and I hope that this matches your requirements.
Thin Film Interference in Oil: Thin film interference can be seen in this oil slick.
How it Works
shows a diagram of how thin film interference works. As light strikes the surface of a film it is either transmitted or reflected at the
upper surface. Light that is transmitted reaches the bottom surface and may once again be transmitted or reflected. The light
reflected from the upper and lower surfaces will interfere. The degree of constructive or destructive interference between the two
light waves is dependent upon the difference in their phase. This difference is dependent upon the thickness of the film layer, the
refractive index of the film, and the angle of incidence of the original wave on the film. Additionally, a phase shift of 180° or ππ
radians may be introduced upon reflection at a boundary depending on the refractive indices of the materials on either side said
boundary. This phase shift occurs if the refractive index of the medium the light is travelling through is less than the refractive
index of the material it is striking. In other words, if n1<n2n1<n2and the light is travelling from material 1 to material 2, then a
phase shift will occur upon reflection. The pattern of light that results from this interference can appear either as light and dark
bands or as colorful bands depending upon the source of the incident light.
26.3.2 https://phys.libretexts.org/@go/page/16196
Light on a Thin Film: Light incident on a thin film. Demonstration of the optical path length difference for light reflected from the
upper and lower boundaries.
Interference will be constructive if the optical path difference is equal to an integer multiple of the wavelength of light:
2 n2 d cos(θ2 ) = mλ (26.3.1)
where mm is the integer, dd is the thickness of the film, and λλ is the wavelength of light. However, this condition may change if
phase shifts occur upon reflection.
learning objectives
Discuss polarization of electromagnetic waves
Definition of Polarization
As we discusses in previous atoms, light waves are a type of electromagnetic waves, in the visible spectrum. These electromagnetic
(EM) waves are transverse waves. Figure 1 demonstrates that a transverse wave is one oscillates perpendicular to the direction of
the energy transfer. If the wave is traveling from left to right, it is oscillating up and down. Polarization is the property of waves
that allow them to oscillate in more than one direction, but that direction is relative to that of the direction the wave is traveling in.
For an EM wave, the direction of polarization is the direction parallel to the electric field. In Figure 2 you can see that the EM and
magnetic fields are perpendicular to the path of travel. Since the direction of polarization is parallel to the electric field, you can
consider the blue arrows to be the direction of polarization.
Figure 2: An EM wave, such as light, is a transverse wave. The electric and magnetic fields are perpendicular to the direction of
propagation.
26.3.3 https://phys.libretexts.org/@go/page/16196
polarized), and the other, b, horizontally. What happens to these waves as they pass through the polarizer?When wave a, the
vertically oscillating wave is passed through the vertically polarized slot, nothing happens. The wave passes through untouched or
manipulated. When wave b, the horizontally oscillating wave is passed through, the vertically polarized slot blocks the wave, and it
is not passed through at all. Now that we understand the concept of polarization, and how it works, how can we apply this to make
it useful? Look at Figure 4. The image on the left is full of glare, which makes it hard to see what we are looking at. The image on
the right was taken with a polarized lens, so you only see the image, and none of the annoying glare. How this works is diagramed
in Figure 5. Many light sources are unpolarized, and are comprised of many waves in all possible directions. Polarized lenses only
allow one direction of light to pass through, minimizing the unwanted aspects of the light rays, such as glare. Simply, passing light
through a polarized material changes the intensity of the light.
Figure 5: A polarizing filter has a polarization axis that acts as a slit passing through electric fields parallel to its direction. The
direction of polarization of an EM wave is defined to be the direction of its electric field.
Figure 4: These two photographs of a river show the effect of a polarizing filter in reducing glare in light reflected from the surface
of water. Part (b) of this figure was taken with a polarizing filter and part (a) was not. As a result, the reflection of clouds and sky
observed in part (a) is not observed in part (b). Polarizing sunglasses are particularly useful on snow and water. (credit: Amithshs,
Wikimedia Commons)
26.3.4 https://phys.libretexts.org/@go/page/16196
I = I0 cos 2θ (26.3.2)
If you pass light through two polarizing filters, you will get varied effects of polarization. If the two filters are oriented exactly
perpendicular to each other, no light will pass through at all. If they are exactly parallel to each other, there will be no additional
affect from the additional filter.
learning objectives
Calculate angle of reflection of complete polarization from indices of refraction
Polarization by Reflection
In the previous atom we discussed how polarized lenses work. In the case of polarized sunglasses, for example, when you look
through them, reflected light is not entirely filtered out; reflected light can be slightly polarized by the reflection process (as shown
in ). Most light sources produce unpolarized light. When light hits a reflective surface, the vertically polarized aspects of that light
are refracted at that surface. The reflected light is more horizontally polarized. To better remember this, we can think of light as an
arrow and the reflective surface as a target. If the arrow hits the target perpendicularly (vertically polarized), it is going to stick in
the target (be refracted into the surface). If the arrow hits the target on its side (horizontally polarized) then it will bounce right off
(be reflected).
Polarization by Reflection: Unpolarized light has equal amounts of vertical and horizontal polarization. After interaction with a
surface, the vertical components are preferentially absorbed or refracted, leaving the reflected light more horizontally polarized.
This is akin to arrows striking on their sides bouncing off, whereas arrows striking on their tips go into the surface.
Since the light is split into two, and part of it is refracted, the amount of polarization to the reflected light depends on the index of
refraction of the reflective surface. We can use the following equation to determine the angle of reflection at which light will be
completely polarized:
n2
tan θb = (26.3.3)
n1
where: θb = angle of reflection of complete polarization (also known as Brewster’s angle); n1 = index of refraction of medium in
which reflected light will travel; and n2 = index of refraction of medium by which light is reflected.
Polarization by Scattering
Just as unpolarized light can be partially polarized by reflecting, it can also be polarized by scattering (also known as Rayleigh
scattering; illustrated in ). Since light waves are electromagnetic (EM) waves (and EM waves are transverse waves) they will
vibrate the electrons of air molecules perpendicular to the direction in which they are traveling. The electrons then produce
26.3.5 https://phys.libretexts.org/@go/page/16196
radiation (acting like small antennae) that is polarized perpendicular to the direction of the ray. The light parallel to the original ray
has no polarization. The light perpendicular to the original ray is completely polarized. In all other directions, the light scattered by
air will be partially polarized.
Polarization by Scattering: Also known as Rayleigh scattering. Unpolarized light scattering from air molecules shakes their
electrons perpendicular to the direction of the original ray. The scattered light therefore has a polarization perpendicular to the
original direction and none parallel to the original direction.
learning objectives
Describe wave-particle relationship that leads to Rayleigh scattering and apply it to explain common phenomena
Rayleigh Scattering
Rayleigh scattering is the elastic scattering of waves by particles that are much smaller than the wavelengths of those waves. The
particles that scatter the light also need to have a refractive index close to 1. This law applies to all electromagnetic radiation, but in
this atom we are going to focus specifically on why the atmosphere scatters the visible spectrum of electromagnetic waves, also
known as visible light. In this case, the light is scattered by the gas molecules of the atmosphere, and the refractive index of air is 1.
Rayleigh scattering is due to the polarizability of an individual molecule. This polarity describes how much the electric charges in
the molecule will vibrate in an electric field. The formula to calculate the intensity of the scattering for a single particle is as
follows:
4 2
8π α 2
I = I0 (1 + cos θ) (26.3.4)
4 2
λ R
where I is the resulting intensity, I0 is the original intensity, α is the polarizability, λ is the wavelength, R is the distance to the
particle, and θ is the scattering angle.
While you will probably not need to use this formula, it is important to understand that scattering has a strong dependence on
wavelength. From the formula, we can see that a shorter wavelength will be scattered more strongly than a longer one.( The longer
the wavelength, the larger the denominator, and from algebra we know that a larger denominator in a fraction means a smaller
number. )
26.3.6 https://phys.libretexts.org/@go/page/16196
Why are Sunsets Colorful?
shows a sunset. We know why the sky is blue, but why are there all those colors in a sunset? The reddening that occurs near the
horizon is because the light has to pass through a significantly higher volume of air than when the sun is high in the sky. This
increases the Rayleigh scattering effect and removes all blue light from the direct path of the observer. The remaining unscattered
light is of longer wavelengths and so appears orange.
learning objectives
Describe process of dispersion
Dispersion
Dispersion is the spreading of white light into its full spectrum of wavelengths. How does this happen? The index of refraction is
different for every medium that light travels through, as we learned in previous atoms. When a light ray enters a medium with a
different index of refraction, the light is dispersed, as shown in with a prism. When white light enters the prism, it spreads. Since
the index of refraction varies with wavelength, the light refracts at different angles as it exits, causing the exiting light rays to
appear as a rainbow, or as a sequence of decreasing wavelengths, from red to violet.
26.3.7 https://phys.libretexts.org/@go/page/16196
Light and a Glass Prism: (a) A pure wavelength of light falls onto a prism and is refracted at both surfaces. (b) White light is
dispersed by the prism (shown exaggerated). Since the index of refraction varies with wavelength, the angles of refraction vary
with wavelength. A sequence of red to violet is produced, because the index of refraction increases steadily with decreasing
wavelength.
This same principle can be applied to rainbows. Refer to. Rainbows are not only caused by refraction, like prisms, but also
reflection. Light enters a drop of water and is reflected from the back of the droplet. The light is refracted once as it enters the drop,
and again as is exits the drop. In water, the refractive index varies with wavelength, so the light is dispersed.
Light and a Water Droplet: Part of the light falling on this water drop enters and is reflected from the back of the drop. This light
is refracted and dispersed both as it enters and as it leaves the drop.
26.3.8 https://phys.libretexts.org/@go/page/16196
Key Points
When the incident light reaches the thin film, it is partially reflected and passed through to the bottom layer of the film. The
difference in refractive indices of the air and film causes the light to change direction and interfere with the reflected portion of
the ray when it emerges.
This interference can be constructive, producing bright colors, or destructive, producing darker colors.
Interference will be constructive if the optical path difference is equal to an integer multiple of the wavelength of light.
The direction of polarization is parallel to the direction of the electric field associated with an electromagnetic wave, which
includes light waves.
By employing the polarization properties of waves, companies like Polariod how been able to produce materials that filter out
unwanted waves of light, and minimize light intensity.
Polarization materials can be oriented at different angles to produce different effects. The new intensity of light after passing
through these materials can be found using the following formula: I=I0cos2θ.
When unpolarized light hits a reflective surface, the vertically polarized aspects are refracted into the surface. The horizontally
polarized aspects are reflected off the surface and the light is now perceived as partially polarized.
When light is reflected, there is an angle at which this light is completely polarized. This is called Brewster’s angle, after the
Scottish physicst who discovered the law.
Unpolarized light can also become polarized when it is scattered in air (also known as Rayleigh scattering). This occurs due to
the fact that the EM waves cause the electron in air to vibrate, producing radiation and causing polarization of the light.
The phenomenon of light scattering is called Rayleigh scattering; it can happen to any electromagnetic waves. It only occurs
when waves encounter particles that are much smaller than the wavelengths of the waves.
The amount that light is scattered is inversely proportional to the fourth power of the light wavelength. For this reason, shorter-
wavelength light like greens and blues scatter more easily than longer wavelengths like yellows and reds.
As you look closer to the sky’s light source, the sun, the light is scattered less and less because the angle between the sun and
the scattering particles approaches 90 degrees. This is why the sun has a yellowish color when we look at it from Earth while
the rest of the sky appears blue.
In outer space, where there is no atmosphere and therefore no particles to scatter the light, the sky appears black and the sun
appears white.
During a sunset, the light must pass through an increased volume of air. This increases the scattering effect, causing the light in
the direct path of the observer to appear orange rather than blue.
Dispersion is a side effect of the law of refraction. As refraction angles depend on wavelength, when light enters a medium with
a different index of refraction, it can be dispersed, like a prism.
The dispersion of white light can often cause the refracted light to be observed is in order of either increasing or decreasing
wavelength, causing a rainbow effect.
As you can see from the visible spectrum, there are some colors that the brain perceives that are not included. This is because
some colors are a mixture of different wavelengths, like pink and magenta.
Key Terms
incident ray: The ray of light that strikes the surface.
interference: An effect caused by the superposition of two systems of waves, such as a distortion on a broadcast signal due to
atmospheric or other effects.
wavelength: The length of a single cycle of a wave, as measured by the distance between one peak or trough of a wave and the
next; it is often designated in physics as λ, and corresponds to the velocity of the wave divided by its frequency.
oscillate: To swing back and forth, especially if with a regular rhythm.
index of refraction: For a material, the ratio of the speed of light in vacuum to that in the material.
polarization: The production of polarized light; the direction in which the electric field of an electromagnetic wave points.
electromagnetic radiation: radiation (quantized as photons) consisting of oscillating electric and magnetic fields oriented
perpendicularly to each other, moving through space
polarizability: The relative tendency of a system of electric charges to become polarized in the presence of an external electric
field
refraction: Changing of a light ray’s direction when it passes through variations in matter.
reflection: the property of a propagated wave being thrown back from a surface (such as a mirror)
dispersion: The separation of visible light by refraction or diffraction.
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26.4: Applications of Wave Optics
learning objectives
Compare optical and electron microscopy
Microscopes are used to view objects that cannot be seen with the naked eye. In this section we will discuss both optical and
electron microscopy.
Optical Microscopy
You have probably used an optical microscope in a high school science class. In optical microscopy, light reflected from an object
passes through the microscope’s lenses; this magnifies the light. The resultant, magnified image is then seen by the eye. Although
this type of microscopy has many limitations, there are several techniques that use properties of light and optics to enhance the
magnified image:
Bright field: This technique increases the contrast by illuminating the surface on which the objects sit from below.
Oblique illumination: This technique illuminates the object from the side, giving it a three-dimensional appearance and
highlighting features that would otherwise not be visible.
Dark field: This technique is good for improving the contrast of transparent objects. A carefully aligned light source minimizes
the un-scattered light entering the object plane and so only collects the light that is scattered by the object itself.
Dispersion staining: This results in a colored image of a colorless object; it does not actually require that the object be stained.
Phase contrast: This uses the refractive index of an object to show differences in optical density as a difference in contrast.
provides a demonstration of this technique.
Electron Microscopy
Electron microscopes use electron beams to achieve higher resolutions than are possible in optical microscopy. Two kinds of
electron microscopes are:
Transmission electron microscope (TEM): The TEM sends an electron beam through a thin slice of a specimen. The electron
interacts with the specimen and is then transmitted onto photographic paper or a screen. Since electron beams have a much
smaller wavelength than traditional light, the resolution of the resulting image is much higher.
Scanning electron microscope (SEM): The SEM shows details on the surface of a specimen and produces a three-dimensional
view by scanning the specimen. shows an SEM image of pollen.
The Spectrometer
A spectrometer uses properties of light to identify atoms by measuring wavelength and frequency, which are functions of radiated
energy.
learning objectives
Compare design and function of early and modern spectrometers
The Spectrometer
A spectrometer is an instrument used to intensely measure light over a specific portion of the electromagnetic spectrum, to identify
materials. The instrument produces lines, much like those produced from diffraction grating as covered in a previous atom, and
then measures the wavelengths and intensities of those lines.
shows a diagram of how a spectrometer works. The source is placed in front of a mirror, which reflects the light emitted from that
object onto a diffraction grating. This grating then disperses the emitted light to anther mirror which spreads the different resultant
wavelengths and reflects them onto a detector which records the findings. This type of instrument is used in spectroscopy.
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Spectrometer Diagram: This diagram shows the light pathways in a spectrometer.
Spectroscopy
Spectroscopy studies the interaction between matter and radiated energy. This radiated energy is a function of wavelength and
frequency. Every type of atom has its own frequency. When the spectrometer produces a reading, the observer can then use
spectroscopy to identify the atoms and therefore molecules that make up that object.
Spectroscopes
Spectroscopes are used in a variety of fields, such as astronomy and chemistry. They use a diffraction grating, movable slit, and a
photodetector. All of these elements are controlled by a computer, which records the findings. A material is heated to
incandescence and it emits a light that is characteristic of its atomic makeup. Each atom has its own spectroscopic ‘fingerprint’. In
you can see a very simple spectroscope based on a prism. As another example, Sodium produces a double yellow band.
learning objectives
Explain how the Michelson interferometer works
Interferometry
Before we can discuss the Michelson Interferometer, it is important we first understand interferometry—which refers to techniques
that use superimposed waves to extract information about the waves. More simply, it uses the interference these waves experience
to make accurate measurements of the waves. It is used in many areas of science, such as astronomy, engineering, oceanography,
physics, and fiber optics.
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Popular applications of interferometry in industry include the measurement of small displacements, refractive index changes, and
surface irregularities. As shown in previous atoms, when two waves with the same frequency combine, the resulting pattern is
determined by the phase difference between the two. Constructive interference occurs when the waves are in phase, and destructive
interference occurs when they are out of phase. Interferometry uses this principle to combine waves and study the resulting wave in
order to obtain information about the original state of the waves.
Fringes in a Michelson Interferometer: Colored and monochromatic fringes in a Michelson interferometer: (a) White light
fringes where the two beams differ in the number of phase inversions; (b) White light fringes where the two beams have
experienced the same number of phase inversions; and (c) Fringe pattern using monochromatic light (sodium D lines).
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Figure 3: This diagram of a Michelson Interferometer shows the path that the light waves travels in the instrument.
Applications
The Michelson Interferometer has been used for the detection of gravitational waves, as a tunable narrow band filter, and as the
core of Fourier transform spectroscopy. It has played an important role in studies of the upper atmosphere, revealing temperatures
and winds (employing both space-borne and ground-based instruments) by measuring the Doppler widths and shifts in the spectra
of airglow and aurora. The best known application of the Michelson Interferometer is the Michelson-Morley experiment—a failed
attempt to demonstrate the effect of the hypothetical ” aether wind” on the speed of light. Their experiment left theories of light
based on the existence of a luminiferous aether without experimental support, and served ultimately as an inspiration for special
relativity.
LCDs
Liquid crystal displays use liquid crystals which do not emit light, but use the light modulating properties of the crystals.
learning objectives
Explain how liquid crystal displays produce images and discuss their benefits and deficiencies
LCDs
LCD stands for a liquid crystal display. The liquid crystals themselves do not emit light, but the display uses the light modulating
properties of the crystals. LCDs can be used to display arbitrary images, such as in a computer monitor or television, by using a
large number of very small pixels, or they can be used to display fixed images, like a digital clock, such as in.
Digital Clock: A digital clock which uses LCD to either hide or display fixed images.
Unlike the newer cathode ray tube (CRT) and plasma displays, LCDS do not use phosphors. For this reason they do not suffer
image burn-in. They do however suffer image persistence. Image burn-in occurs when an image is displayed so many times, or for
so long, that an outline of image can be seen even when the display is turned off. Image persistence is similar, but the outline fades
away shortly after the display is turned off and is not permanent.
LCD displays are made up of numerous layers. A typical layer is diagrammed in. Each pixel of an LCD consists of a layer of
molecules aligned between two transparent electrodes and two polarizing films, and the actual liquid crystals are between these
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polarizing filters. The light passes through the first filter, and is blocked by the second. The electrodes are used to align the crystals
in a particular direction, which produces the image seen on the screen. The crystals do not emit any light, but rather give the light a
specific shape to be emitted in.
Layers of LCD Displays: Polarizing filter film with a vertical axis to polarize light as it enters.Glass substrate with ITO electrodes.
The shapes of these electrodes will determine the shapes that will appear when the LCD is turned on. Vertical ridges etched on the
surface are smooth.Twisted nematic liquid crystal.Glass substrate with common electrode film (ITO) with horizontal ridges to line
up with the horizontal filter.Polarizing filter film with a horizontal axis to block/pass light.Reflective surface to send light back to
viewer. (In a backlit LCD, this layer is replaced with a light source. )
learning objectives
Explain how information is stored on the optical disks
Overview
Compact disks (CDs) and digital video disks (DVDs) are examples of optical discs. They are read in an optical disc drive which
directs a laser beam at the disc. The reader then detects whether the beam has been reflected or scattered.
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Early Version of an Optical Disc: In this early version of an optical disc, you can see the pits and lands which either reflect back
light or scatter it.
These microscopic pits and lands cover the entire surface of the disc in a spiral path, starting in the center and working its way
outward. The data is stored either by a stamping machine or laser and is read when the data is illuminated by a laser diode in the
disc drive. The disc spins at a faster speed when it is being read in the center track, and slower for an outer track. This is because
the center tracks are smaller in circumference and therefore can be read quicker.
These pits also act as slits and cause the light to be diffracted as it is reflected back, which causes an iridescent effect. This explains
the rainbow pattern that you see on the back of a CD, as shown in.
Compact Disc: The bottom surface of a compact disc showing characteristic iridescence.
Key Points
In optical microscopy, light reflected from an object passes through the microscope’s lenses; this magnifies the light. The
resultant, magnified image is then seen by the eye. This technique has many limitations but can be enhanced in various ways to
create more contrast.
The transmission electron microscope (TEM) sends an electron beam through a thin slice of a specimen. The electron is then
transmitted onto photographic paper or a screen. Since electron beams have a much smaller wavelength than traditional light,
the resolution of this image is much higher.
The scanning electron microscope (SEM) shows details on the surface of a specimen and produces a three-dimensional view by
scanning the specimen.
The source is placed in front of a mirror, which reflects the light emitted from that object onto a diffraction grating. This grating
then disperses the emitted light to anther mirror which spreads the different resultant wavelengths and reflects them onto a
detector which records the findings.
Early forms of spectrometers were simple prisms, but modern spectrometers are automated by a computer and can record a
much broader range of frequencies.
Spectrometers are used in spectroscopy. Spectroscopy studies the interaction between matter and radiated energy. This radiated
energy is a function of wavelength and frequency. Every type of atom has its own frequency.
Interferometry refers to techniques that use superimposed waves to extract information about the waves.
Michelson interferometer works by splitting a beam of light into two paths, bouncing them back and recombining them to
create an interference pattern. To create interference fringes on a detector, the paths may be different lengths or composed of
different materials.
The best known application of the Michelson interferometer is the Michelson-Morley experiment, the unexpected null result of
which was an inspiration for special relativity.
LCDs use an electric field to arrange the liquid crystals into the desired pattern, and then pass light through these layers to
produce an image on the screen.
LCDs can be used to display an arbitrary image made up of tiny fixed pixels, or can be used to display a fixed image, as in on a
digital clock.
A twisted nematic display is the most common LCD in use. This type of display is on calculators, digital watches, and clocks.
When no electric field is applied, the molecules are twisted, and let some light through. When the field is applied, they untwist,
blocking the light and are seen as black.
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The iridescent layer of the disc is imprinted with tiny pits and lands. Pits scatter light when illuminated, and produce a reading
of 0; lands reflect light back and produce a reading of 1.
The optical disc drive records the 0 and 1 readings and translates them into binary data which is used to relay whatever
information is recorded on the disc.
Rainbow pattern on the back of a CD is due diffraction of the reflected light by pits.
Key Terms
microscopy: using microscopes to view objects that cannot be seen with the naked eye
contrast: A difference in lightness, brightness and/or hue between two colors that makes them more or less distinguishable
incandescence: Incandescence is the emission of light (visible electromagnetic radiation) from a hot body as a result of its
temperature.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
superimposed: Positioned on or above something else, especially in layers
interference: An effect caused by the superposition of two systems of waves, such as a distortion on a broadcast signal due to
atmospheric or other effects.
LCD: a liquid crystal display.
helical: In the shape of a helix, twist.
nematic: Describing the structure of some liquid crystals whose molecules align in loose parallel lines.
binary data: Data which can take on only two possible values, traditionally termed 0 and 1.
pit: An imprint on an optical disc that scatters light when illuminated.
land: A flat area on an optical disc that reflects light when illuminated.
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CHAPTER OVERVIEW
27: Special Relativity is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
27.1: Introduction
learning objectives
Explain why the Galilean invariance didn’t work in Maxwell’s equations
Galilean invariance or Galilean relativity states that the laws of motion are the same in all inertial (or non-accelerating) frames.
Galileo Galilei first described this principle in 1632 using the example of a ship travelling at constant velocity, without rocking, on
a smooth sea; any observer doing experiments below the deck would not be able to tell whether the ship was moving or stationary.
The fact that the Earth orbits around the sun at approximately 30 km/s offers a somewhat more dramatic example, though it is
technically not an inertial reference frame.
Specifically, the term Galilean invariance today usually refers to this principle as applied to Newtonian mechanics—that is,
Newton’s laws hold in all inertial frames. In this context it is sometimes called Newtonian relativity. Among the axioms from
Newton’s theory are:
There exists an absolute space in which Newton’s laws are true. An inertial frame is a reference frame in relative uniform
motion to absolute space.
All inertial frames share a universal (or absolute) time.
Derivation
Galilean relativity can be shown as follows. Consider two inertial frames S and S’. A physical event in S will have position
coordinates r = (x, y, z) and time t; similarly for S’. By the second axiom above, one can synchronize the clock in the two frames
and assume t = t‘. Suppose S’ is in relative uniform motion to S with velocity v. Consider a point object whose position is given by
r = r(t) in S. We see that
′
r (t) = r(t) − vt (27.1.1)
This transformation of variables between two inertial frames is called Galilean transformation. Now, the velocity of the particle is
given by the time derivative of the position:
Galilean Invariance: Newtonian mechanics is invariant under a Galilean transformation between observation frames (shown).
This is called Galilean invariance.
′
d ′
d
u (t) = r (t) = r(t) − v = u(t) − v (27.1.2)
dt dt
′
d ′
d
a (t) = u (t) = u(t) − 0 = a(t) (27.1.3)
dt dt
It is this simple but crucial result that implies Galilean relativity. Assuming that mass is invariant in all inertial frames, the above
equation shows that Newton’s laws of mechanics, if valid in one frame, must hold for all frames. But it is assumed to hold in
absolute space, therefore Galilean relativity holds.
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Issues
Both Newtonian mechanics and the Maxwell’s equations were well established by the end of the 19th century. The puzzle lied in
the fact that the Galilean invariance didn’t work in Maxwell’s equations. That is, unlike Newtonian mechanics, Maxwell’s
equations are not invariant under a Galilean transformation. Albert Einstein’s central insight in formulating special relativity was
that, for full consistency with electromagnetism, mechanics must also be revised, such that Lorentz invariance (introduced later)
replaces Galilean invariance. At the low relative velocities characteristic of everyday life, Lorentz invariance and Galilean
invariance are nearly the same, but for relative velocities close to that of light they are very different.
Einstein’s Postulates
Special relativity is based on Einstein’s two postulates: the Principle of Relativity and the Principle of Invariant Light Speed.
learning objectives
Identify two postulates forming the foundation of special relativity
In the late 19th century, the Newtonian mechanics was considered to be valid in all inertial frames of reference, which are moving
at a constant relative velocity with respect to each other. (See our previous lesson on “Galilean-Newtonian Relativity. “) One issue,
however, was that another well-established theory, the laws of electricity and magnetism represented by Maxwell’s equations, was
not “invariant” under Galilean transformation—meaning that Maxwell’s equations don’t maintain the same forms for different
inertial frames. In his “Special Theory of Relativity,” Einstein resolved the puzzle and broadened the scope of the invariance to
extend the validity of all physical laws, including electromagnetic theory, to all inertial frames of reference.
Albert Einstein: Albert Einstein, a true pioneer of modern physics. His work on relativity, gavity, quantum mechanics, and
statistical physics revolutionized physics.
Einstein’s Postulates
With two deceptively simple postulates and a careful consideration of how measurements are made, Einstein produced the theory
of special relativity.
1. The Principle of Relativity: The laws of physics are the same and can be stated in their simplest form in all inertial frames of
reference.
This postulate relates to reference frames. It says that there is no preferred frame and, therefore, no absolute motion.
2. The Principle of Invariant Light Speed: The speed of light c is a constant, independent of the relative motion of the source and
observer.
The laws of electricity and magnetism predict that light travels at c = 2.998×108 m/s in a vacuum, but they do not specify the frame
of reference in which light has this speed. Physicists assumed that there exists a stationary medium for the propagation of light,
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which they called ” luminiferous aether. ” In 1887, Michelson and Morley attempted to detect the relative motion of the Earth
through the stationary luminiferous aether, but their negative results implied the speed of light c is independent of the motion of the
source relative to the observer. Einstein accepted the result of the experiment and incorporated it in his theory of relativity.
This postulate might sound easy to accept, but it is rather counterintuitive. Imagine that you can throw a baseball at a speed v
(relative to you). If you are on a train moving at a speed V and throw a ball in the direction of the train’s movement, the baseball
will travel at a speed v+V for an observer stationary on the ground.
Now, instead of a baseball, let’s say you have a laser pointer. You turn on the laser pointer while you are on a moving train. What
would be the speed of light from the laser pointer for a stationary observer on the ground? Our intuition says that it should be c+V.
However, Einstein says that it should be only c!
learning objectives
Discuss the invariance of the speed of light and identify the value of that speed in vacuum
The speed of light in vacuum, commonly denoted c, is a universal physical constant that is crucial to many areas of physics. Its
value is 299,792,458 m/s; this is a precise known value because the length of the meter is itself derived from this constant and the
international standard for time. According to special relativity, c is the maximum speed at which all energy, matter, and information
in the universe can travel. It is the speed at which all massless particles and associated fields (including electromagnetic radiation
such as light) travel in vacuum. It is also the speed of gravity (i.e., of gravitational waves) predicted by current theories. Such
particles and waves (including light) travel at c regardless of the motion of the source or the inertial frame of reference of the
observer. In the theory of relativity, c interrelates space and time in the Lorentz transformation; it also appears in the famous
equation of mass-energy equivalence: E = mc2.
Sunlight’s Flight to Earth: Sunlight takes about 8 minutes and 19 seconds to reach the earth (based on the average distance
between the sun and the earth)
First Measurement
The first quantitative estimate of the speed of light was made in 1676 by Rømer. From the observation that the periods of Jupiter’s
innermost moon (Io) appeared to be shorter when Earth was approaching Jupiter than when it was moving away, he concluded that
light travels at a finite speed. He estimated that it takes light 22 minutes to cross the diameter of Earth’s orbit. Christiaan Huygens
combined this estimate with an estimate for the diameter of the Earth’s orbit to obtain an estimate of the speed of light of 220,000
km/s, 26 percent lower than the actual value.
Key Points
Galilean invariance states that Newton’s laws hold in all inertial frames.
Newtonian mechanics assumes that there exists an absolute space and that time is universal.
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Albert Einstein’s central insight in formulating special relativity was that, for full consistency with electromagnetism,
mechanics must also be revised such that Lorentz invariance (introduced later) replaces Galilean invariance.
The special theory of relativity explores the consequences of the invariance of c with the assumption that the laws of physics are
the same in all inertial frames of reference.
Einstein’s first postulate says that the laws of physics are the same and can be stated in their simplest form in all inertial frames
of reference. It means that there is no preferred frame and, therefore, no absolute motion.
The speed of light c is a constant, independent of the relative motion of the source and observer.
The speed at which light waves propagate in vacuum is independent both of the motion of the wave source and of the inertial
frame of reference of the observer. This invariance of the speed of light was postulated by Einstein in 1905 in his work on
special relativity.
The value of the speed of light is 299,792,458 m/s; this is a precise known value because the length of the meter is itself derived
from this constant and the international standard for time.
c is the maximum speed at which all energy, matter, and information in the universe can travel.
Key Terms
Lorentz invariance: First introduced by Lorentz in an effort to explain how the speed of light was observed to be independent
of the reference frame, and to understand the symmetries of the laws of electromagnetism.
absolute space: A concept introduced by Newton that assumes space remains always similar and immovable.
Maxwell’s equations: A set of equations describing how electric and magnetic fields are generated and altered by each other
and by charges and currents.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
luminiferous aether: Light-bearing aether; the postulated medium for the propagation of light.
Lorentz transformation: a transformation relating the spacetime coordinates of one frame of reference to another in special
relativity
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27.2: Consequences of Special Relativity
learning objectives
Formulate conclusions of the theory of special relativity, noting the assumptions that were made in deriving it
The relativity of simultaneity is the concept that simultaneity–whether two events occur at the same time–is not absolute, but
depends on the observer’s frame of reference.
According to the theory of special relativity, it is impossible to say in an absolute sense whether two distinct events occur at the
same time if those events are separated in space, such as a car crash in London and another in New York. The question of whether
the events are simultaneous is relative: in some reference frames the two accidents may happen at the same time, in other frames
(in a different state of motion relative to the events) the crash in London may occur first, and still in other frames, the New York
crash may occur first. If the two events are causally connected (“event A causes event B”), then the relativity of simultaneity
preserves the causal order (i.e. “event A causes event B” in all frames of reference).
If we imagine one reference frame assigns precisely the same time to two events that are at different points in space, a reference
frame that is moving relative to the first will generally assign different times to the two events. This is illustrated in the ladder
paradox, a thought experiment which uses the example of a ladder moving at high speed through a garage.
A mathematical form of the relativity of simultaneity (“local time”) was introduced by Hendrik Lorentz in 1892, and physically
interpreted (to first order in v/c) as the result of a synchronization using light signals by Henri Poincaré in 1900. However, both
Lorentz and Poincaré based their conceptions on the aether as a preferred but undetectable frame of reference, and continued to
distinguish between “true time” (in the aether) and “apparent” times for moving observers.
In 1905, Albert Einstein abandoned the (classical) aether and emphasized the significance of relativity of simultaneity to our
understanding of space and time. He deduced the failure of absolute simultaneity from two stated assumptions: 1) the principle of
relativity–the equivalence of inertial frames, such that the laws of physics apply equally in all inertial coordinate systems; 2) the
constancy of the speed of light detected in empty space, independent of the relative motion of its source.
Observer Standing on the Platform: Reference frame of an observer standing on the platform (length contraction not depicted).
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Observer Onboard the Train: The train-and-platform experiment from the reference frame of an observer onboard the train.
Time Dilation
Time dilation is an actual difference of elapsed time between two events as measured by observers moving relative to each other.
learning objectives
Explain why time dilation can be ignored in daily life
Time dilation is an actual difference of elapsed time between two events as measured by observers either moving relative to each
other.
For instance, two rocket ships (A and B) speeding past one another in space would experience time dilation. If they somehow had a
clear view into each other’s ships, each crew would see the others’ clocks and movement as going too slowly. That is, inside the
frame of reference of Ship A, everything is moving normally, but everything over on Ship B appears to be moving slower (and vice
versa).
From a local perspective, time registered by clocks that are at rest with respect to the local frame of reference (and far from any
gravitational mass) always appears to pass at the same rate. In other words, if a new ship, Ship C, travels alongside Ship A, it is “at
rest” relative to Ship A. From the point of view of Ship A, new Ship C’s time would appear normal too.
The formula for determining time dilation is: Δt ′
= γΔt =
Δt
2 2
√1−v / c
where Δt is the time interval between two co-local events (i.e. happening at the same place) for an observer in some inertial frame
(e.g. ticks on his clock), this is known as the proper time, Δt’ is the time interval between those same events, as measured by
another observer, inertially moving with velocity v with respect to the former observer, v is the relative velocity between the
observer and the moving clock, c is the speed of light, and γ = 1
√1−v2 / c2
is the Lorentz factor. Thus the duration of the clock cycle of a moving clock is found to be increased: it is measured to be “running
slow”. Note that for speeds below 1/10 the speed of light, Lorentz factor is approximately 1. Thus, time dilation effects and
extremely small and can be safely ignored in a daily life. They become important only when an object approaches speeds on the
order of 30,000 km/s (1/10 the speed of light).
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Lorentz Factor: Lorentz factor as a function of speed (in natural units where c = 1). Notice that for small speeds (less than 0.1), γ
is approximately 1.
Effects of Time Dilation: The Twin Paradox and the Decay of the Muon
The twin paradox is a thought experiment: one twin makes a journey into space and returns home to find that twin remained aged
more.
learning objectives
Explain the twin paradox within the standard framework of special relativity
The twin paradox is a thought experiment in special relativity involving identical twins, one of whom makes a journey into space in
a high-speed rocket and returns home to find that the twin who remained on Earth has aged more.
This occurs because special relativity shows that the faster one travels, the slower time moves for them.
This result appears puzzling because each twin sees the other twin as traveling, and so, according to a naive application of time
dilation, each should paradoxically find the other to have aged more slowly. In other words, from the perspective of the rocketship,
the earth is traveling away from the ship and from the perspective of the earth, the rocket is traveling away.
However, this scenario can be resolved within the standard framework of special relativity (because the twins are not equivalent;
the space twin experienced additional, asymmetrical acceleration when switching direction to return home), and therefore is not a
paradox in the sense of a logical contradiction.
The Earth and the ship are not in a symmetrical relationship: regardless of whether we view the situation from the perspective of
the Earth or the ship, the ship experiences additional acceleration forces. The ship has a turnaround in which it accelerates and
changes direction whereas the earth does not. Since there is no symmetry, it is not paradoxical if one twin is younger than the other.
Nevertheless twin paradox is useful as a demonstration that special relativity is self-consistent.
In the spacetime diagram, drawn for the reference frame of the Earth-based twin, that twin’s world line coincides with the vertical
axis (his position is constant in space, moving only in time). On the first leg of the trip, the second twin moves to the right (black
sloped line); and on the second leg, back to the left. Blue lines show the planes of simultaneity for the traveling twin during the first
leg of the journey; red lines, during the second leg. Just before turnaround, the traveling twin calculates the age of the Earth-based
twin by measuring the interval along the vertical axis from the origin to the upper blue line. Just after turnaround, if he recalculates,
he’ll measure the interval from the origin to the lower red line. In a sense, during the U-turn the plane of simultaneity jumps from
blue to red and very quickly sweeps over a large segment of the world line of the Earth-based twin. The traveling twin reckons that
there has been a jump discontinuity in the age of the Earth-based twin.
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Spacetime Diagram of the Twin Paradox: Spacetime diagram of the twin paradox. Time is relative, but both twins are not
equivalent (the ship experiences additional acceleration to changes the direction of travel).
Length Contraction
Objects that are moving undergo a length contraction along the dimension of motion; this effect is only significant at relativistic
speeds.
learning objectives
Explain why length contraction can be ignored in daily life
Length contraction is the physical phenomenon of a decrease in length detected by an observer of objects that travel at any non-
zero velocity relative to that observer. Length contraction arises due to the fact that the speed of light in a vacuum is constant in any
frame of reference. By taking this into account, as well as some geometrical considerations, we will show how perceived time and
length are affected.
Example 27.2.1:
Let us imagine a simple clock system that consists of two mirrors A and B in a vacuum. A light pulse bounces between the two
mirrors. The separation of the mirrors is L, and the clock ticks once each time the light pulse hits a given mirror. Now imagine
that the clock is at rest. The time that it will take for the light pulse to go from mirror A to mirror B and then back to mirror A
be can be described by:
2L
Δt = (27.2.1)
c
where c is the speed of light. Now imagine that the clock is moving in the horizontal direction relative to a stationary observer.
The light pulse is emitted from mirror A. To the stationary observer, it appears that the light pulse has a longer path to travel
because by the time the light reaches mirror B the clock has already moved somewhat in the horizontal direction. This is the
same case for the light pulse on its way back. The stationary observer will perceive that it will take the light a total of:
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Geometry for a Clock at Rest: This illustrates the path that light must traverse when the clock is at rest.
2D
′
Δt = (27.2.2)
c
to traverse its path. We can see that D is longer than L, so that means that.
Geometry for a Moving Clock: This illustrates the path that light must traverse when the clock is moving from the
perspective of a stationary observer.
Example 27.2.2:
We have established that in a frame of reference that is moving relative to the clock (the stationary observer is moving in the
clock’s frame of reference), the clock appears to run more slowly. Now let us imagine that we want to measure the length of a
ruler. This time let us imagine that you are moving with velocity v. You can mathematically determine the length of the ruler in
your frame of reference (L‘) by multiplying your velocity (v) by the time that you perceive that it takes you to pass by the ruler
(t‘). Expressing this in equation form, L‘ = vt‘. Now, if someone in the ruler’s rest frame wanted to determine the length of the
ruler, they could do the following. They could mathematically determine the length of the ruler in their frame of reference (L)
by multiplying your velocity (v) by the time that they perceive that it takes you to pass by the ruler (t). This is expressed in the
following equation: L = vt. Just as in the clock explanation, the ruler appears to be moving in your frame of reference, so t will
be longer than t‘ (your time interval). Consequently, the length of the ruler will appear to be shorter in your frame of reference
(the phenomenon of length contraction occurred).
The effect of length contraction is negligible at everyday speeds and can be ignored for all regular purposes. Length
contraction becomes noticeable at a substantial fraction of the speed of light (as illustrated in ) with the contraction only in the
direction parallel to the direction in which the observed body is travelling.
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Observed Length of an Object: Observed length of an object at rest and at different speeds
Example 27.2.3:
For example, at a speed of 13,400,000 m/s (30 million mph,.0447c), the length is 99.9 percent of the length at rest; at a speed
of 42,300,000 m/s (95 million mph, 0.141c), the length is still 99 percent. As the magnitude of the velocity approaches the
speed of light, the effect becomes dominant.The mathematical formula for length contraction is:
−−−−−−
2
L0 v
L = = L0 √ 1 − (27.2.3)
2
γ(v) c
where L0 is the proper length (the length of the object in its rest frame); L is the length observed by an observer in relative
motion with respect to the object; v is the relative velocity between the observer and the moving object; c is the speed of light;
and the Lorentz factor is defined as:
1
γ(v) = (27.2.4)
− −−−− −−−
2 2
√ 1 − v /c
In this equation it is assumed that the object is parallel with its line of movement. For the observer in relative movement, the
length of the object is measured by subtracting the simultaneously measured distances of both ends of the object. An observer
at rest viewing an object traveling very close to the speed of light would observe the length of the object in the direction of
motion as very close to zero.
Key Points
According to the theory of special relativity, it is impossible to say in an absolute sense whether two distinct events occur at the
same time if those events are separated in space.
A mathematical form of the relativity of simultaneity was introduced by Hendrik Lorentz and physically interpreted by Henri
Poincaré. The conceptions were based on the aether as a preferred but undetectable frame of reference.
Albert Einstein deduced the failure of absolute simultaneity from two assumptions: 1) the principle of relativity; 2) the
constancy of the speed of light detected in empty space, independent of the relative motion of its source.
Time dilation effects are extremely small for speeds below 1/10 the speed of light and can be safely ignored at daily life.
Time dilation effects become important when an object approaches speeds on the order of 30,000 km/s (1/10 the speed of light).
The formula for determining time dilation is: Δt =′ Δt
.
2 2
√1−v / c
From a naive application of time dilation, each twin should paradoxically find the other to have aged more slowly.
The scenario is resolved within the standard framework of special relativity: the twins are not equivalent, the space twin
experiences additional, asymmetrical acceleration when switching direction to return home.
Twin paradox is useful as a demonstration that special relativity is self-consistent.
Length contraction is negligible at everyday speeds and can be ignored for all regular purposes.
27.2.6 https://phys.libretexts.org/@go/page/16200
Length contraction becomes noticeable at a substantial fraction of the speed of light with the contraction only in the direction
parallel to the direction in which the observed body is travelling.
An observer at rest viewing an object travelling very close to the speed of light would observe the length of the object in the
direction of motion as very near zero.
Key Terms
aether: A space-filling substance or field, thought to be necessary as a transmission medium for the propagation of
electromagnetic or gravitational forces.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
speed of light: the speed of electromagnetic radiation in a perfect vacuum: exactly 299,792,458 meters per second by definition
time dilation: The slowing of the passage of time experienced by objects in motion relative to an observer; measurable only at
relativistic speeds.
Lorentz factor: The factor, used in special relativity, to calculate the degree of time dilation, length contraction and relativistic
mass of an object moving relative to an observer.
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27.3: Relativistic Quantities
learning objectives
Express velocity-addition formulas for objects at speeds much less and approaching the speed of light
A velocity-addition formula is an equation that relates the velocities of moving objects in different reference frames.
As Galileo Galilei observed in 17th century, if a ship is moving relative to the shore at velocity vv, and a fly is moving with
velocity uu as measured on the ship, calculating the velocity of the fly as measured on the shore is what is meant by the addition of
the velocities vv and uu. When both the fly and the ship are moving slowly compared to speed of light, it is accurate enough to use
the vector sum s=u+vs=u+v where ss is the velocity of the fly relative to the shore.
According to the theory of special relativity, the frame of the ship has a different clock rate and distance measure, and the notion of
simultaneity in the direction of motion is altered, so the addition law for velocities is changed.
Since special relativity dictates that the speed of light is the same in all frames of reference, light shone from the front of a moving
car can’t go faster than light from a stationary lamp. Since this is counter to what Galileo used to add velocities, there needs to be a
new velocity addition law.
This change isn’t noticeable at low velocities but as the velocity increases towards the speed of light it becomes important. The
addition law is also called a composition law for velocities. For collinear motions, the velocity of the fly relative to the shore is
given by the following equation:
v+u
s = (27.3.1)
2
1 + vu/c
Composition law for velocities gave the first test of the kinematics of the special theory of relativity. Using a Michelson
interferometer, Hyppolite Fizeau measured the speed of light in a fluid moving parallel to the light in 1851. The speed of light in
the fluid is slower than the speed of light in vacuum, and it changes if the fluid is moving along with the light. The speed of light in
a collinear moving fluid is predicted accurately by the collinear case of the relativistic formula.
Setup of the Fizeau Experiment: A light ray emanating from the source S’ is reflected by a beam splitter G and is collimated into
a parallel beam by lens L. After passing the slits O1 and O2, two rays of light travel through the tubes A1 and A2, through which
water is streaming back and forth as shown by the arrows. The rays reflect off a mirror m at the focus of lens L’, so that one ray
always propagates in the same direction as the water stream, and the other ray opposite to the direction of the water stream. After
passing back and forth through the tubes, both rays unite at S, where they produce interference fringes that can be visualized
through the illustrated eyepiece. The interference pattern can be analyzed to determine the speed of light traveling along each leg of
the tube.
Relativistic Momentum
Relativistic momentum is given as γ m 0v where m is the object’s invariant mass and γ is Lorentz transformation.
0
learning objectives
Compare Newtonian and relativistic momenta for objects at speeds much less and approaching the speed of light
Relativistic Momentum
Newtonian physics assumes that absolute time and space exist outside of any observer. This gives rise to Galilean relativity, which
states that the laws of motion are the same in all inertial frames. It also results in a prediction that the speed of light can vary from
one reference frame to another. However, this is contrary to observation. In the theory of special relativity, Albert Einstein keeps
the postulate that the equations of motion do not depend on the reference frame, but assumes that the speed of light c is invariant.
As a result, position and time in two reference frames are related by the Lorentz transformation instead of the Galilean
transformation.
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Albert Einstein: Albert Einstein in 1921
Consider, for example, a reference frame moving relative to another at velocity v in the x direction. The Galilean transformation
gives the coordinates of the moving frame as
′
t =t (27.3.2)
′
x = x − vt (27.3.3)
′
x = γ(x − vt) (27.3.5)
Conservation laws in physics, such as the law of conservation of momentum, must be invariant. That is, the property that needs to
be conserved should remain unchanged regardless of changes in the conditions of measurement. This means that the conservation
law needs to hold in any frame of reference. Newton’s second law [with mass fixed in the expression for momentum (p=m*v)], is
not invariant under a Lorentz transformation. However, it can be made invariant by making the inertial mass m of an object a
function of velocity:
m = γm0 (27.3.7)
It is important to note that for speeds much less than the speed of light, Newtonian momentum and relativistic momentum are
approximately the same. As one approaches the speed of light, however, relativistic momentum becomes infinite while Newtonian
momentum continues to increases linearly. Thus, it is necessary to employ the expression for relativistic momentum when one is
dealing with speeds near the speed of light.
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Relativistic and Newtonian Momentum: This figure illustrates that relativistic momentum approaches infinity as the speed of
light is approached. Newtonian momentum increases linearly with speed.
learning objectives
Evaluate possibility for an object to travel at the speed of light
Relativistic and Newtonian Kinetic Energy: This figure illustrates how relativistic and Newtonian Kinetic Energy are related to
the speed of an object. The relativistic kinetic energy increases to infinity when an object approaches the speed of light, this
indicates that no body with mass can reach the speed of light. On the other hand, Newtonian kinetic energy continues to increase
without bound as the speed of an object increases.
Relativistic Mass
Relativistic mass was defined by Richard C. Tolman pictured left of Albert Einstein here in 1934 as which holds for all particles,
including those moving at the speed of light. For a slower than light particle, a particle with a nonzero rest mass, the formula
27.3.3 https://phys.libretexts.org/@go/page/16203
becomes where is the rest mass and is the Lorentz factor. The Lorentz factor is equal to: γ =
1
, where v is the relative
2 2
√1−v / c
Richard C. Tolman and Albert Einstein: Richard C. Tolman (1881 – 1948) with Albert Einstein (1879 – 1955) at Caltech, 1932
When the relative velocity is zero, is simply equal to 1, and the relativistic mass is reduced to the rest mass. As the velocity
increases toward the speed of light (c), the denominator of the right side approaches zero, and consequently approaches infinity.
In the formula for momentum the mass that occurs is the relativistic mass. In other words, the relativistic mass is the
proportionality constant between the velocity and the momentum.
While Newton’s second law remains valid in the form the derived form is not valid because in is generally not a constant.
Relativistic Energy
−−−−−−−−−−−− √( E
2 2
−(pc)
2
Relativistic energy (Er
2
= √(m0 c )
2
+ (pc) ) is connected with rest mass via the following equation: m = c
2
. Here the
term represents the square of the Euclidean norm (total vector length) of the various momentum vectors in the system, which
reduces to the square of the simple momentum magnitude, if only a single particle is considered. This equation reduces to when the
momentum term is zero. For photons where the equation reduces to.
Today, the predictions of relativistic energy and mass are routinely confirmed from the experimental data of particle accelerators
such as the Relativistic Heavy Ion Collider. The increase of relativistic momentum and energy is not only precisely measured but
also necessary to understand the behavior of cyclotrons and synchrotron, which accelerate particles to near the speed of light.
learning objectives
Describe properties of antiparticles
Antimatter is material composed of antiparticles, which have the same mass as particles of ordinary matter but have opposite
charge and quantum spin. Antiparticles bind with each other to form antimatter in the same way that normal particles bind to form
normal matter. For example, a positron (the antiparticle of the electron, with symbol e+) and an antiproton (symbol p–) can form an
antihydrogen atom. Furthermore, mixing matter and antimatter can lead to the annihilation of both, in the same way that mixing
antiparticles and particles does. This gives rise to high-energy photons (gamma rays) and other particle-antiparticle pairs. The end
result of antimatter meeting matter is a release of energy proportional to the mass, as shown in the mass-energy equivalence
equation, E = mc2.
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Antihydrogen and Hydrogen Atoms: Antihydrogen consists of an antiproton and a positron; hydrogen consists of a proton and an
electron.
Almost all matter observable from the earth seems to be made of matter rather than antimatter. If antimatter-dominated regions of
space existed, the gamma rays produced in annihilation reactions along the boundary between matter and antimatter regions would
be detectable.
Antimatter may still exist in relatively large amounts in far-away galaxies due to cosmic inflation in the primordial time of the
universe. Antimatter galaxies, if they exist, are expected to have the same chemistry and absorption and emission spectra as
normal-matter galaxies, and their astronomical objects would be observationally identical, making them difficult to distinguish
from normal-matter galaxies.
There is considerable speculation as to why the observable universe is apparently composed almost entirely of matter (as opposed
to a mixture of matter and antimatter), whether there exist other places that are almost entirely composed of antimatter instead, and
what sorts of technology might be possible if antimatter could be harnessed. At this time, the apparent asymmetry of matter and
antimatter in the visible universe is one of the greatest unsolved problems in physics.
2
− mc
2
where mm is rest mass, v is velocity, c is speed of light.
√1−(v/c) )
learning objectives
Compare classical and relativistic kinetic energies for objects at speeds much less and approaching the speed of light
In classical mechanics, the kinetic energy of an object depends on the mass of a body as well as its speed. The kinetic energy is
equal to the mass multiplied by the square of the speed, multiplied by the constant 1/2. The equation is given as:
1
2
Ek = mv (27.3.10)
2
where mm is the mass and vv is the speed (or the velocity) of the body.
The classical kinetic energy of an object is related to its momentum by the equation:
2
p
Ek = (27.3.11)
2m
where pp is momentum.
If the speed of a body is a significant fraction of of the speed of light, it is necessary to employ special relativity to calculate its
kinetic energy. It is important to know how to apply special relativity to problems with high speed particles. In special relativity, we
must change the expression for linear momentum. Using mm for rest mass, v and V for the object’s velocity and speed
respectively, and cc for the speed of light in vacuum, the relativistic expression for linear momentum is:
27.3.5 https://phys.libretexts.org/@go/page/16203
Time Magazine – July 1, 1946: The popular connection between Einstein, E = mc , and the atomic bomb was prominently
2
indicated on the cover of Time magazine (July 1946) by the writing of the equation on the mushroom cloud itself.
p = mγv , where γ is the Lorentz factor:
1
γ = −−−−−−− −
2
√1 − (v/c)
Since the kinetic energy of an object is related to its momentum, we intuitively know that the relativistic expression for kinetic
energy will also be different from its classical counterpart. Indeed, the relativistic expression for kinetic energy is:
2
mc 2
Ek = − −−−− −−−− − mc (27.3.12)
2
√ 1 − (v/c) )
The equation shows that the energy of an object approaches infinity as the velocity vv approaches the speed of light cc. Thus it is
impossible to accelerate an object across this boundary.
The mathematical by-product of this calculation is the mass-energy equivalence formula (referred to in ). The body at rest must
have energy content equal to:
2
Erest = E0 = mc (27.3.13)
The general expression for the kinetic energy of an object that is not at rest is:
KE = mc
2 2
− m0 c , where m is the relativistic mass of the object and m0 is the rest mass of the object.
At a low speed (v<<cv<<c), the relativistic kinetic energy may be approximated well by the classical kinetic energy. We can show
this to be true by using a Taylor expansion for the reciprocal square root and keeping first two terms of the relativistic kinetic
energy equation. When we do this we get:
2
1 2 2 2
1 2
Ek ≈ mc (1 + v / c ) − mc = mv (27.3.14)
2 2
Thus, the total energy can be partitioned into the energy of the rest mass plus the traditional classical kinetic energy at low speeds.
Key Points
When two objects are moving slowly compared to speed of light, it is accurate enough to use the vector sum of velocities:
s = u+v .
As the velocity increases towards the speed of light, the vector sum of velocities is replaced with: s = . v+u
1+vu/c2
Composition law for velocities gave the first test of the kinematics of the special theory of relativity when, using a Michelson
interferometer, Hyppolite Fizeau measured the speed of light in a fluid moving parallel to the light.
27.3.6 https://phys.libretexts.org/@go/page/16203
Newtonian physics assumes that absolute time and space exist outside of any observer, resulting in a prediction that the speed of
light can vary from one reference frame to another.
In the theory of special relativity, the equations of motion do not depend on the reference frame while the speed of light (c) is
invariant.
Within the domain of classical mechanics, relativistic momentum closely approximates Newtonian momentum. At low velocity
γm0vγm0v is approximately equal to m0vm0v, the Newtonian expression for momentum.
In special relativity, an object that has a mass cannot travel at the speed of light.
Relativistic mass is defined as m =rel
E
2
c
and can be viewed as the proportionality constant between the velocity and the
momentum.
−−−−−−−−−−−−
Relativistic energy is connected with rest mass via the following equation: E r
2
= √(m0 c )
2
+ (pc)
2
.
The end result of antimatter meeting matter is a release of energy proportional to the mass, as shown in the mass-energy
equivalence equation,E = mc . 2
Although almost all matter observable from the Earth seems to be made of matter rather than antimatter, antimatter may still
exist in relatively large amounts in far-away galaxies.
No explanation for the apparent asymmetry of matter and antimatter in the visible universe exists.
Relativistic kinetic energy equation shows that the energy of an object approaches infinity as the velocity approaches the speed
of light. Thus it is impossible to accelerate an object across this boundary.
Kinetic energy calculations lead to the mass -energy equivalence formula: E = E = mc .
rest 0
2
At a low speed (v << C ), the relativistic kinetic energy may be approximated by the classical kinetic energy. Thus, the total
energy can be partitioned into the energy of the rest mass plus the traditional Newtonian kinetic energy at low speeds.
Key Terms
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
interferometer: Any of several instruments that use the interference of waves to determine wavelengths and wave velocities,
determine refractive indices, and measure small distances, temperature changes, stresses, and many other useful measurements.
speed of light: the speed of electromagnetic radiation in a perfect vacuum: exactly 299,792,458 meters per second by definition
Galilean transformation: a transformation used to transform between the coordinates of two reference frames which differ
only by constant relative motion within the constructs of Newtonian physics.
Lorentz transformation: a transformation relating the spacetime coordinates of one frame of reference to another in special
relativity
Lorentz factor: The factor, used in special relativity, to calculate the degree of time dilation, length contraction and relativistic
mass of an object moving relative to an observer.
rest mass: the mass of a body when it is not moving relative to an observer
annihilation: the process of a particle and its corresponding antiparticle combining to produce energy
positron: The antimatter equivalent of an electron, having the same mass but a positive charge.
antimatter: matter that is composed of the antiparticles of those that constitute normal matter
classical mechanics: All of the physical laws of nature that account for the behaviour of the normal world, but break down
when dealing with the very small (see quantum mechanics) or the very fast or very heavy (see relativity).
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27.4: Implications of Special Relativity
learning objectives
Formulate major changes in the understanding of time, space, mass, and energy that were introduced by the theory of
Special Relativity
The theory of Special Relativity and its implications spurred a paradigm shift in our understanding of the nature of the universe, the
fundamental fabric of which being space and time. Before 1905, scientists considered space and time as completely independent
objects. Time could not affect space and space could not affect time. After 1905, however, the Special Theory of Relativity
destroyed this old, but intuitive, view. Specifically, Special Relativity showed us that space and time are not independent of one
another but can be mixed into each other and therefore must be considered as the same object, which we shall denote as space-time.
The consequences of space/time mixing are:
time dilation
and length contraction.
Other important consequences which will be discussed in another section are
Relativity of Simultaneity (for certain events, the sequence in which they occur is dependent on the observer)
Nothing can move faster than the speed of light (we shall denote the value of the speed of light as cc)
Why is there a mixing between space and time? In order to examine this we must know the founding principles of relativity. They
are:
1. The Principle of Relativity: The laws of physics for observers which are not accelerating relative to one another should be the
same.
2. The Principle of Invariant Light Speed: All observers, moving at constant speed, measure the same speed of light regardless of
how fast they are moving.
If one accepts the second principle as fact then it immediately follows that space and time are not independent. Why? Let us look at
the experiment in in which there is a light source, a fixed observer, and a rocket ship moving toward the light source. The two
observers are related by a coordinate , or space-time, transformation. The second principle then tells us that no matter how fast the
rocket is moving, both observers must measure the same light speed emanating from the source. The only way this can happen is if
the clock of the rocket observer is ticking at a different rate than the stationary observer. How much different? This can be
expressed in the time dilation equation:
Measuring Light: A stationary observer will measure the same speed of light as an observer who is moving in a rocket ship even if
that rocket is moving close to light speed.
tm
ts = (27.4.1)
−−−−−−−−
2
v
√1 − ( )
c2
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Where tsts is the time elapsed for the stationary observer, tmtm is the time elapsed for the moving observer, and vv is the velocity
of the rocket as measured from the stationary frame. One can then see that as V → C then the time elapsed in the stationary frame
goes to infinity. To place some numbers, let v = 0.99986c, and t = 1 sec . The time elapsed in the stationary frame is then one
m
hour. Thus for every second lived in the rocket, the stationary man lives one hour!
One of the more radical results of time dilation is the so-called “twin paradox. ” The twin paradox is a thought experiment in
special relativity that involves identical twins. One of the twins goes on a journey into space in a rocket that has a velocity near the
speed of light. Upon returning home the twin finds that the twin that remained on Earth has aged more. This consequence altered
the perception that aging is necessarily constant.
The square root factor in the time dilation equation is very important and we denote it as:
This factor shows up frequently in special relativity. For example the length contraction formula is:
−−−−−− −
v L0
L = L0 √ 1 − ( ) = (27.4.2)
c γ
where L is the rest length, the length of an object measured in the co-moving frame of the object, and L is the length of the object
0
as measured by the observer who sees the object moving at speed v . In our rocket example, the stationary observer measures the
length of the rocket as being less than what someone who was moving with the rocket would measure. This altered the perception
that the length of an object would appear the same regardless of the reference frame of the observer.
Another radical finding that was made possible by the discovery of special relativity is the equivalence of energy and mass.
Combined with other laws of physics, the postulates of special relativity predict that mass and energy are related by: E = mc , 2
where c is the speed of light in vacuum. This altered the perception that mass and energy are completely different things and paved
the way for nuclear power and nuclear weapons.
As a final note it is important to note how big the speed of light is compared to every day life. The speed of light is:
8
c = 3 × 10 m/s (27.4.3)
Thus in every day life γ ≈ 1 and we do not experience significant time dilation or length contraction. If we did, life would be very
different.
Four-Dimensional Space-Time
We live in four-dimensional space-time, in which the ordering of certain events can depend on the observer.
learning objectives
Formulate major results of special relativity
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Two Coordinate Systems: Two coordinate systems in which the primed frame moves with velocity v with respect to the unprimed
frame
You should not find it odd to work with four dimensions; any time you have to meet your friend somewhere you have to tell him
four variables: where (three spatial coordinates) and when (one time coordinate). In other words, we have always lived in four
dimensions, but so far you have probably thought of space and time as completely separate.
where, for example, Δt = t − t0 ; t is the time at which the measurement took place; and t0 is the time at which the light was
turned on.
Similarly, observer A’ measures:
′ ′ ′ ′
(Δt , Δx , Δy , Δz )
where we set up the system such that both observers agree on (t0 , x0 , y , z0 )
0
. Due to the invariance of the speed of light both
observers will agree on:
2 2 2
ΔX + ΔY + ΔZ 2 2 2 2 2 2
=c → 0 = −c ΔT + ΔX + ΔY + ΔZ (27.4.4)
2
ΔT
where (T, X, Y, Z)refers to the coordinates in either frame. Therefore there is a specific rule that all light paths must follow. For
general events we can define the quantity:
2 2 2 2 2 2
s = −c Δt + Δx + Δy + Δz (27.4.5)
This is known as the line element and is the same for all observers. (Using the principles of relativity, you can prove this for general
separations, not just light rays). The set of all coordinate transformations that leave the above quantity invariant are known as
Lorentz Transformations. It follows that the coordinate systems of all physical observers are related to each other by Lorentz
Transformations. (The set of all Lorentz transformations form what mathematicians call a group, and the study of group theory has
revolutionized physics). For the coordinate transformation in, the transformations are:
′ 2
t = γ (t − vx/ c )
′
x = γ(x − vt)
(27.4.6)
′
y =y
′
z =z
where
1
γ = − −−−−−− − (27.4.7)
2
√ 1 − (v/c)
2
s < 0 : time-like (27.4.9)
2
s = 0 : null (27.4.10)
We separate these events because they are very different. For example, for a space-like separation you can always find a coordinate
transformation that reverses the time-ordering of the events (try to prove this for the example in ). This phenomenon is known as
the relativity of simultaneity and may be counterintuitive.
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Space-Like Space-Time Separations
Let us examine two car crashes, one in New York and one in London, such that they occur in the same time in one frame. In this
situation, the space-time separation between the two events is space-like. The question of whether the events are simultaneous is
relative: in some reference frames the two accidents may happen at the same time; in other frames (in a different state of motion
relative to the events) the crash in London may occur first; and in still other frames the New York crash may occur first. (In
practice, this does not affect us because the frames that switch the order of the events must move at unreachably high speeds. )
Special Relativity
Finally, let’s discuss an important result of special relativity — that the energy EE of an object moving with speed vv is:
2
m0 c 2
E = − −−−−−− − = mc (27.4.11)
2
√ 1 − (v/c)
where m0m0 is the mass of the object at rest and m=γm0m=γm0 is the mass when the object is moving. The above formula
immediately shows why it is impossible to travel faster than the speed of light. As v → c, m → ∞ , and it takes an infinite amount
of energy to accelerate the object further.
learning objectives
Identify factors that affect motion near massive objects
Relativity
Special relativity indicates that humans live in a four-dimensional space-time where the ‘distance’ ssbetween points in space-time
can be regarded as:
\[\mathrm{ s } ^ { 2 } = - \mathrm { c } ^ { 2 } \Delta \mathrm { t } ^ { 2 } + \Delta \mathrm { x } ^ { 2 } + \Delta \mathrm{ y } ^
{ 2 } + \Delta \mathrm { z } ^ { 2 } = \mathrm{ X } ^ { \mathrm { T } } \eta \mathrm{ X }
The last equation is a matrix relation in which TT denotes transpose (for vectors A B = A ⋅ B ), X is the vector \)( \mathrm { c }
T
\Delta \mathrm { t } , \Delta \mathrm { x } , \Delta \mathrm { y } , \Delta \mathrm { z } )\), and η is the matrix:
−1 0 0 0
⎛ ⎞
⎜ 0 1 0 0⎟
η =⎜ ⎟ (27.4.12)
⎜ 0 0 1 0⎟
⎝ ⎠
0 0 0 1
A matrix that goes in between two vectors to give a length is called a metric. In mathematics, a metric or distance function is a
function which defines a distance between elements of a set. In this case, the set is the space-time and the elements are points in
that space-time. A space-time with the ηη metric is called Minkowski space and ηη is the Minkowski metric.
Four-dimensional Minkowski space-time is only one of many different possible space-times (geometries) which differ in their
metric matrix. An arbitrary metric matrix can be denoted as gg, raising questions as to what space-times with different metrics
represent. In 1916, Einstein found the importance of these space-times in his theory of general relativity.
General Relativity
General relativity, or the general theory of relativity, is the geometric theory of gravitation published by Albert Einstein in 1916. It
is the current description of gravitation in modern physics. General relativity generalizes special relativity and Newton’s law of
universal gravitation, providing a unified description of gravity as a geometric property of space and time, or space-time. In
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particular, the curvature of space-time is directly related to the energy and momentum of whatever matter and radiation are present.
The relation is specified by the Einstein field equations, a system of partial differential equations.
g ↔ curvature ↔ Energy and Momentum (27.4.13)
People can use the metric to calculate curvature and then use the Einstein field equations to relate the curvature to the energy and
momentum of the space-time. Going in the reverse order, energy and momentum affect the curvature and the space-time. Thus,
energy and momentum curves space-time. Minkowski space is the special space devoid of matter, and as a result, it is completely
flat. The precise definition of curvature requires knowledge of advanced mathematics, but an intuitive way to understand it is that
the definition of a straight line changes in curved spacetime.
A pictorial example is shown in the figure below, where anything that is near the Earth has its motion altered due to the curved
space-time. This altering of motion affects satellites, the moon, and even humans. If the space around the Earth were flat, humans
would simply float away if they jumped upwards. Since the Earth alters the space-time, humans are pulled toward the Earth. In
essence, gravity is a geometrical effect.
Key Points
Time is relative.
Lengths of moving objects are different than if their lengths were measured in a co-moving frame.
Time and space are not independent.
We live in a four-dimensional universe; the first three dimensions are spatial, and the fourth is time.
The coordinate system of physical observers are related to one another via Lorentz transformations.
Nothing can travel faster than the speed of light.
The metric matrix can be used to calculate the curvature of space-time.
Curvature can be related to energy and momentum via the Einstein equations.
Due to the curvature of space-time, the motion of objects are altered near massive objects.
Key Terms
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
time dilation: The slowing of the passage of time experienced by objects in motion relative to an observer; measurable only at
relativistic speeds.
length contraction: Observers measure a moving object’s length as being smaller than it would be if it were stationary.
line element: An invariant quantity in special relativity
Lorentz transformation: a transformation relating the spacetime coordinates of one frame of reference to another in special
relativity
relativity of simultaneity: For space-like separated space-time points, the time-ordering between events is relative.
general relativity: A theory extending special relativity and uniformly accounting for gravity and accelerated frames of
reference, postulating that space-time curves in the presence of mass.
Minkowski space: A four dimensional flat space-time. Because it is flat, it is devoid of matter.
metric: A metric, or distance function, is a function which defines a distance between elements of a set.
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CHAPTER OVERVIEW
28: Introduction to Quantum Physics is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
28.1: History and Quantum Mechanical Quantities
learning objectives
Explain how the photoelectric effect paradox was solved by Albert Einstein.
Electrons are emitted from matter when light shines on a surface . This is called the photoelectric effect, and the electrons emitted
in this manner are called photoelectrons.
The Photoelectric Effect: Electrons are emitted from matter by absorbed light.
The photoelectric effect typically requires photons with energies from a few electronvolts to 1 MeV for heavier elements, roughly
in the ultraviolet and X-ray range. Study of the photoelectric effect led to important steps in understanding the quantum nature of
light and electrons and influenced the formation of the concept of wave-particle duality. The photoelectric effect is also widely used
to investigate electron energy levels in matter.
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Photoelectric Effect: A brief introduction to the Photoelectric Effect and electron photoemission.
Heinrich Hertz discovered the photoelectric effect in 1887. Although electrons had not been discovered yet, Hertz observed that
electric currents were produced when ultraviolet light was shined on a metal. By the beginning of the 20th century, physicists
confirmed that:
The energy of the individual photoelectrons increased with the frequency (or color) of the light, but was independent of the
intensity (or brightness) of the radiation.
The photoelectric current was determined by the light’s intensity; doubling the intensity of the light doubled the number of
emitted electrons.
This observation was very puzzling to many physicists. At the time, light was accepted as a wave phenomenon. Since energy
carried by a wave should only depend on its amplitude (and not on the frequency of the wave), the frequency dependence of the
emitted electrons’ energies didn’t make sense.
In 1905, Albert Einstein solved this apparent paradox by describing light as composed of discrete quanta (now called photons),
rather than continuous waves. Building on Max Planck’s theory of black body radiation, Einstein theorized that the energy in each
quantum of light was equal to the frequency multiplied by a constant h, later called Planck’s constant. A photon above a threshold
frequency has the required energy to eject a single electron, creating the observed effect. As the frequency of the incoming light
increases, each photon carries more energy, hence increasing the energy of each outgoing photoelectron. By doubling the number
of photons as the intensity is doubled, the number of photelectrons should double accordingly.
According to Einstein, the maximum kinetic energy of an ejected electron is given by K max = hf − ϕ , where h is the Planck
constant and f is the frequency of the incident photon. The term ϕ is known as the work function, the minimum energy required to
remove an electron from the surface of the metal. The work function satisfies ϕ = hf , where f is the threshold frequency for the
0 0
metal for the onset of the photoelectric effect. The value of work function is an intrinsic property of matter.
Is light then composed of particles or waves? Young’s experiment suggested that it was a wave, but the photoelectric effect
indicated that it should be made of particles. This question would be resolved by de Broglie: light, and all matter, have both wave-
like and particle-like properties.
learning objectives
Compare photon energy with the frequency of the radiation
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Electromagnetic spectrum: This shows the electromagnetic spectrum, including the visible region, as a function of both frequency
(left) and wavelength (right).
Maxwell’s equations predicted an infinite number of frequencies of electromagnetic waves, all traveling at the speed of light. This
was the first indication of the existence of the entire electromagnetic spectrum. Maxwell’s predicted waves included waves at very
low frequencies compared to infrared, which in theory might be created by oscillating charges in an ordinary electrical circuit of a
certain type. In 1886, the physicist Hertz built an apparatus to generate and detect what are now called radio waves, in an attempt to
prove Maxwell’s equations and detect such low-frequency electromagnetic radiation. Hertz found the waves and was able to infer
(by measuring their wavelength and multiplying it by their frequency) that they traveled at the speed of light. Hertz also
demonstrated that the new radiation could be both reflected and refracted by various dielectric media, in the same manner as light.
λ
E
h
hc
λ
where νν is the frequency, λ is the wavelength, E is photon energy, c is the speed of light,
and hh is the Planck constant. Generally, electromagnetic radiation is classified by wavelength into radio waves, microwaves,
terahertz (or sub-millimeter) radiation, infrared, the visible region humans perceive as light, ultraviolet, X-rays, and gamma rays.
The behavior of EM radiation depends on its wavelength. When EM radiation interacts with single atoms and molecules, its
behavior also depends on the amount of energy per quantum (photon) it carries.
Most parts of the electromagnetic spectrum are used in science for spectroscopic and other probing interactions as ways to study
and characterize matter. Also, radiation from various parts of the spectrum has many other uses in communications and
manufacturing.
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learning objectives
State physical properties of a photon
A photon is an elementary particle, the quantum of light. It has no rest mass and has no electric charge. The modern photon concept
was developed gradually by Albert Einstein to explain experimental observations of the photoelectric effect, which did not fit the
classical wave model of light. In particular, the photon model accounted for the frequency dependence of light’s energy. Max
Planck explained black body radiation using semiclassical models, in which light is still described by Maxwell’s equations, but the
material objects that emit and absorb light, do so in amounts of energy that are quantized.
Photons are emitted in many natural processes. They are emitted from light sources such as floor lamps or lasers . For example,
when a charge is accelerated it emits photons, a phenomenon known as synchrotron radiation. During a molecular, atomic or
nuclear transition to a lower or higher energy level, photons of various energy will be emitted or absorbed respectively. A photon
can also be emitted when a particle and its corresponding antiparticle are annihilated. During all these processes, photons will carry
energy and momentum.
E = pc . Combining this with Eq. 1, we get p = . Here, λ is the wavelength of the light. Since momentum is a vector
hν h
=
c λ
quantity and p points in the direction of the photon’s propagation, we can write p = hk , where h = h
2π
and is k a wave vector.
You may wonder how an object with zero rest mass can have nonzero momentum. This confusion often arises because of the
commonly used form of momentum (mvmv in non-relativistic mechanics and γmvγmv in relativistic mechanics, where vv is
velocity and γ = . ) This formula, obviously, shouldn’t be used in the case v = c .
1
2
v
√1−
2
c
learning objectives
Explain importance of quantum mechanics for technology and other branches of science
The field of quantum mechanics has been enormously successful in explaining many of the features of our world. The behavior of
the subatomic particles (electrons, protons, neutrons, photons, and others) that make up all forms of matter can often be
satisfactorily described only using quantum mechanics. Quantum mechanics has also strongly influenced string theory.
Quantum mechanics is also critically important for understanding how individual atoms combine covalently to form molecules.
The application of quantum mechanics to chemistry is known as quantum chemistry. Relativistic quantum mechanics can, in
principle, mathematically describe most of chemistry. Quantum mechanics can also provide quantitative insight into ionic and
covalent bonding processes by explicitly showing which molecules are energetically favorable to which other molecules and the
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magnitudes of the energies involved. Furthermore, most of the calculations performed in modern computational chemistry rely on
quantum mechanics.
A great number of modern technological inventions operate on a scale where quantum effects are significant. Examples include the
laser , the transistor (and thus the microchip), the electron microscope, and magnetic resonance imaging (MRI). The study of
semiconductors led to the invention of the diode and the transistor, which are indispensable parts of modern electronic systems and
devices.
Particle-Wave Duality
Wave–particle duality postulates that all physical entities exhibit both wave and particle properties.
learning objectives
Describe experiments that demonstrated wave-particle duality of physical entities
Wave–particle duality postulates that all physical entities exhibit both wave and particle properties. As a central concept of
quantum mechanics, this duality addresses the inability of classical concepts like “particle” and “wave” to fully describe the
behavior of (usually) microscopic objects.
From a classical physics point of view, particles and waves are distinct concepts. They are mutually exclusive, in the sense that a
particle doesn’t exhibit wave-like properties and vice versa. Intuitively, a baseball doesn’t disappear via destructive interference,
and our voice cannot be localized in space. Why then is it that physicists believe in wave-particle duality? Because that’s how
mother Nature operates, as they have learned from several ground-breaking experiments. Here is a short, chronological list of those
experiments:
Young’s double-slit experiment: In the early Nineteenth century, the double-slit experiments by Young and Fresnel provided
evidence that light is a wave. In 1861, James Clerk Maxwell explained light as the propagation of electromagnetic waves
according to the Maxwell’s equations.
Black body radiation: In 1901, to explain the observed spectrum of light emitted by a glowing object, Max Planck assumed that
the energy of the radiation in the cavity was quantized, contradicting the established belief that electromagnetic radiation is a
wave.
Photoelectric effect: Classical wave theory of light also fails to explain photoelectric effect. In 1905, Albert Einstein explained
the photoelectric effects by postulating the existence of photons, quanta of light energy with particulate qualities.
De Broglie’s wave (matter wave): In 1924, Louis-Victor de Broglie formulated the de Broglie hypothesis, claiming that all
matter, not just light, has a wave-like nature. His hypothesis was soon confirmed with the observation that electrons (matter)
also displays diffraction patterns, which is intuitively a wave property.
From these historic achievements, physicists now accept that all entities in nature behave as both a particle and a wave, depending
on the specifics of the phenomena under consideration. Because of its counter-intuitive aspect, the meaning of the particle-wave
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duality is still a point of debate in quantum physics. The standard interpretation is that the act of measurement causes the set of
probabilities, governed by a probability distribution function acquired from a “wave”, to immediately and randomly assume one of
the possible values, leading to a “particle”-like result.
So, why do we not notice a baseball acting like a wave? The wavelength of the matter wave associated with a baseball, say moving
at 95 miles per hour, is extremely small compared to the size of the ball so that wave-like behavior is never noticeable.
Diffraction Revisited
De Broglie’s hypothesis was that particles should show wave-like properties such as diffraction or interference.
learning objectives
Compare application of X-ray, electron, and neutron diffraction for materials research
The de Broglie hypothesis, formulated in 1924, predicts that particles should also behave as waves. The wavelength of an electron
is given by the de Broglie equation λ = . Here h is Planck’s constant and pp the relativistic momentum of the electron. λ is
h
Electron Diffraction Pattern: Typical electron diffraction pattern obtained in a transmission electron microscope with a parallel
electron beam.
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Neutrons have also been used for studying crystalline structures. They are scattered by the nuclei of the atoms, unlike X-rays,
which are scattered by the electrons of the atoms. Thus, neutron diffraction has some key differences compared to more common
methods using X-rays or electrons. For example, the scattering of X-rays is highly dependent on the atomic number of the atoms
(i.e., the number of electrons), whereas neutron scattering depends on the properties of the nuclei. In addition, the magnetic
moment of the neutron is non-zero, and can thus also be scattered by magnetic fields. This means that neutron scattering is more
useful for determining the properties of atomic nuclei, despite the fact that neutrons are significantly harder to create, manipulate,
and detect compared to X-rays and electrons.
learning objectives
Relate the wave function with the probability density of finding a particle, commenting on the constraints the wave function
must satisfy for this to make sense
In quantum mechanics, a wave function is a probability amplitude describing the quantum state of a particle and how it behaves.
Typically, its values are complex numbers. For a single particle, it is a function of space and time. The most common symbols for a
wave function are ψ(x) or Ψ(x) (lowercase or uppercase psi, respectively), when the wave function is given as a function of
position x . Although ψ is a complex number, |ψ|2 is a real number and corresponds to the probability density of finding a particle
in a given place at a given time, if the particle’s position is measured.
Trajectories of a Harmonic Oscillator: This figure shows some trajectories of a harmonic oscillator (a ball attached to a spring) in
classical mechanics (A-B) and quantum mechanics (C-H). In quantum mechanics (C-H), the ball has a wave function, which is
shown with its real part in blue and its imaginary part in red. The trajectories C-F are examples of standing waves, or “stationary
states. ” Each standing-wave frequency is proportional to a possible energy level of the oscillator. This “energy quantization” does
not occur in classical physics, where the oscillator can have any energy.
The laws of quantum mechanics (the Schrödinger equation) describe how the wave function evolves over time. The wave function
behaves qualitatively like other waves, such as water waves or waves on a string, because the Schrödinger equation is
mathematically a type of wave equation. This explains the name “wave function” and gives rise to wave-particle duality.
The wave function must satisfy the following constraints for the calculations and physical interpretation to make sense:
It must everywhere be finite.
It must everywhere be a continuous function and continuously differentiable.
It must everywhere satisfy the relevant normalization condition so that the particle (or system of particles) exists somewhere
with 100-percent certainty.
If these requirements are not met, it’s not possible to interpret the wave function as a probability amplitude. This is because the
values of the wave function and its first order derivatives may not be finite and definite (having exactly one value), which means
that the probabilities can be infinite and have multiple values at any one position and time, which is nonsense. Furthermore, when
we use the wave function to calculate an observation of the quantum system without meeting these requirements, there will not be
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finite or definite values to use (in this case the observation can take a number of values and can be infinite). This is not a possible
occurrence in a real-world experiment. Therefore, a wave function is meaningful only if these conditions are satisfied.
learning objectives
Formulate the de Broglie relation as an equation
In quantum mechanics, the concept of matter waves (or de Broglie waves) reflects the wave-particle duality of matter. The theory
was proposed by Louis de Broglie in 1924 in his PhD thesis. The de Broglie relations show that the wavelength is inversely
proportional to the momentum of a particle, and is also called de Broglie wavelength.
Einstein derived in his theory of special relativity that the energy and momentum of a photon has the following relationship:
E = pc (E: energy, p: momentum, c : speed of light).
He also demonstrated, in his study of photoelectric effects, that energy of a photon is directly proportional to its frequency, giving
us this equation:
E = hν (h: Planck constant, ν : frequency).
Combining the two equations, we can derive a relationship between the momentum and wavelength of light:
p=
E
c
=
hν
c
=
h
λ
. Therefore, we arrive at λ = h
p
.
De Broglie’s hypothesis is that this relationship λ = hp , derived for electromagnetic waves, can be adopted to describe matter (e.g.
electron, neutron, etc.) as well.
De Broglie didn’t have any experimental proof at the time of his proposal. It took three years for Clinton Davisson and Lester
Germer to observe diffraction patterns from electrons passing a crystalline metallic target (see ). Before the acceptance of the de
Broglie hypothesis, diffraction was a property thought to be exhibited by waves only. Therefore, the presence of any diffraction
effects by matter demonstrated the wave-like nature of matter. This was a pivotal result in the development of quantum mechanics.
Just as the photoelectric effect demonstrated the particle nature of light, the Davisson–Germer experiment showed the wave-nature
of matter, thus completing the theory of wave-particle duality.
Davisson-Germer Experimental Setup: The experiment included an electron gun consisting of a heated filament that released
thermally excited electrons, which were then accelerated through a potential difference (giving them a certain amount of kinetic
energy towards the nickel crystal). To avoid collisions of the electrons with other molecules on their way towards the surface, the
experiment was conducted in a vacuum chamber. To measure the number of electrons that were scattered at different angles, an
electron detector that could be moved on an arc path about the crystal was used. The detector was designed to accept only
elastically scattered electrons.
Experiments with Fresnel diffraction and specular reflection of neutral atoms confirm the application to atoms of the de Broglie
hypothesis. Further, recent experiments confirm the relations for molecules and even macromolecules, normally considered too
large to undergo quantum mechanical effects. In 1999, a research team in Vienna demonstrated diffraction for molecules as large as
fullerenes. The researchers calculated a De Broglie wavelength of the most probable C velocity as 2.5 pm.
60
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The Heisenberg Uncertainty Principle
The uncertainty principle asserts a basic limit to the precision with which some physical properties of a particle can be known
simultaneously.
learning objectives
Relate the Heisenberg uncertainty principle with the matter wave nature of all quantum objects
The uncertainty principle is any of a variety of mathematical inequalities, asserting a fundamental limit to the precision with which
certain pairs of physical properties of a particle, such as position x and momentum pp or energy E and time t , can be known
simultaneously. The more precisely the position of some particle is determined, the less precisely its momentum can be known, and
¯¯
¯
2
, where σ is the standard deviation of position, σp is the
x
¯
¯¯
standard deviation of momentum, and h = . The uncertainty principle is inherent in the properties of all wave-like systems, and
h
2π
it arises in quantum mechanics simply due to the matter wave nature of all quantum objects. Thus, the uncertainty principle actually
states a fundamental property of quantum systems, and is not a statement about the observational success of current technology.
The principle is quite counterintuitive, so the early students of quantum theory had to be reassured that naive measurements to
violate it were bound always to be unworkable. One way in which Heisenberg originally illustrated the intrinsic impossibility of
violating the uncertainty principle is by using an imaginary microscope as a measuring device.
Heisenberg Microscope: Heisenberg’s microscope, with cone of light rays focusing on a particle with angle ϵ He imagines an
experimenter trying to measure the position and momentum of an electron by shooting a photon at it.
Example 28.1.1:
Example One
If the photon has a short wavelength and therefore a large momentum, the position can be measured accurately. But the photon
scatters in a random direction, transferring a large and uncertain amount of momentum to the electron. If the photon has a long
wavelength and low momentum, the collision does not disturb the electron’s momentum very much, but the scattering will
reveal its position only vaguely.
Example Two
If a large aperture is used for the microscope, the electron’s location can be well resolved (see Rayleigh criterion); but by the
principle of conservation of momentum, the transverse momentum of the incoming photon and hence the new momentum of
the electron resolves poorly. If a small aperture is used, the accuracy of both resolutions is the other way around.
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Heisenberg’s Argument
Heisenberg’s argument is summarized as follows. He begins by supposing that an electron is like a classical particle, moving in the
xx direction along a line below the microscope, as in the illustration to the right. Let the cone of light rays leaving the microscope
lens and focusing on the electron makes an angle εε with the electron. Let λ be the wavelength of the light rays. Then, according to
the laws of classical optics, the microscope can only resolve the position of the electron up to an accuracy of δx = When λ
sin(ε/2)
an observer perceives an image of the particle, it’s because the light rays strike the particle and bounce back through the
microscope to their eye. However, we know from experimental evidence that when a photon strikes an electron, the latter has a
recoil with momentum proportional to h/λ, where h is Planck’s constant.
It is at this point that Heisenberg introduces objective indeterminacy into the thought experiment. He writes that “the recoil cannot
be exactly known, since the direction of the scattered photon is undetermined within the bundle of rays entering the microscope”.
In particular, the electron’s momentum in the xx direction is only determined up to δpx≈hλsin(ε/2)δpx≈hλsin(ε/2). Combining the
relations for δx and δpx , we thus have that δx ⋅ δpx ≈ (
λ
sin(ε/2)
)(
h
λ
sin(ε/2)) = h , which is an approximate expression of
Heisenberg’s uncertainty principle.
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Philosophical Implications
Since its inception, many counter-intuitive aspects of quantum mechanics have provoked strong philosophical debates.
learning objectives
Formulate the Copenhagen interpretation of the probabilistic nature of quantum mechanics
Since its inception, many counter-intuitive aspects and results of quantum mechanics have provoked strong philosophical debates
and many interpretations. Even fundamental issues, such as Max Born’s basic rules interpreting ψ ∗ ψ as a probability density
function took decades to be appreciated by society and many leading scientists. Indeed, the renowned physicist Richard Feynman
once said, “I think I can safely say that nobody understands quantum mechanics. ”
Niels Bohr and Albert Einstein: Niels Bohr (left) and Albert Einstein (right). Despite their pioneering contributions to the
inception of the quantum mechanics, they disagreed on its interpretation.
The Copenhagen interpretation has philosophical implications to the concept of determinism. According to the theory of
determinism, for everything that happens there are conditions such that, given those conditions, nothing else could happen.
Determinism and free-will seem to be mutually exclusive. If the universe, and any person in it are governed by strict and universal
laws , then that means that a person’s behavior could be predicted based on sufficient knowledge of the circumstances obtained
prior to that person’s behavior. However, the Copenhagen interpretation suggests a universe in which outcomes are not fully
determined by prior circumstances but also by probability. This gave thinkers alternatives to strictly bound possibilities, proposing
a model for a universe that follows general rules but never had a predetermined future.
Philosophical Implications
It is also believed therein that any well-defined application of the quantum mechanical formalism must always make reference to
the experimental arrangement. This is due to the quantum mechanical principle of wave function collapse. That is, a wave function
which is initially in a superposition of several different possible states appears to reduce to a single one of those states after
interaction with an observer. In simplified terms, it is the reduction of the physical possibilities into a single possibility as seen by
an observer. This raises philosophical questions about whether something that is never observed actually exists.
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Einstein-Podolsky-Rosen (EPR) Paradox
Albert Einstein (shown in , himself one of the founders of quantum theory) disliked this loss of determinism in measurement in the
Copenhagen interpretation. Einstein held that there should be a local hidden variable theory underlying quantum mechanics and,
consequently, that the present theory was incomplete. He produced a series of objections to the theory, the most famous of which
has become known as the Einstein-Podolsky-Rosen (EPR) paradox. John Bell showed by Bell’s theorem that this “EPR” paradox
led to experimentally testable differences between quantum mechanics and local realistic theories. Experiments have been
performed confirming the accuracy of quantum mechanics, thereby demonstrating that the physical world cannot be described by
any local realistic theory. The Bohr-Einstein debates provide a vibrant critique of the Copenhagen Interpretation from an
epistemological point of view.
Quantum Entanglement
One of the most bizarre aspect of the quantum mechanics is known as quantum entanglement. Quantum entanglement occurs when
particles interact physically and then become separated, while isloated from the rest of the universe to prevent any deterioration of
the quantum state. According to the Copenhagen interpretation of quantum mechanics, their shared state is indefinite until
measured. Once a particle in the entangled state is measured and its state is determined, the Copenhagen interpretation demands
that the other particles’ state is also determined instantaneously. This bizarre nature of action at a distance (which seemingly violate
the speed limit on the transmission of information implicit in the theory of relativity) is what bothered Einstein the most.
(According to the theory of relativity, nothing can travel faster than the speed of light in a vacuum. This seemingly puts a limit on
the speed at which information can be transmitted. ) Quantum entanglement is the key element in proposals for quantum computers
and quantum teleportation.
Key Points
The energy of the emitted electrons depends only on the frequency of the incident light, and not on the light intensity.
Einstein explained the photoelectric effect by describing light as composed of discrete particles.
Study of the photoelectric effect led to important steps in understanding the quantum nature of light and electrons, which would
eventually lead to the concept of wave-particle duality.
Electromagnetic radiation is classified according to wavelength, divided into radio waves, microwaves, terahertz (or sub-
millimeter) radiation, infrared, the visible region humans perceive as light, ultraviolet, X-rays, and gamma rays.
Photon energy is proportional to the frequency of the radiation.
Most parts of the electromagnetic spectrum are used in science for spectroscopic and other probing interactions as ways to study
and characterize matter.
E = hv Energy of photon is proportional to its frequency.
Because of different forms of interaction involved, X-ray, electron, and neutron are suitable for different studies of material
properties.
De Broglie’s idea completed the wave-particle duality.
|ψ|2(x) corresponds to the probability density of finding a particle in a given location x at a given time.
The laws of quantum mechanics (the Schrödinger equation) describe how the wave function evolves over time. The
Schrödinger equation is a type of wave equation, which explains the name “wave function”.
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A wave function must satisfy a set of mathematical constraints for the calculations and physical interpretation to make sense.
de Broglie relations show that the wavelength is inversely proportional to the momentum of a particle.
The Davisson-Germer experiment demonstrated the wave-nature of matter and completed the theory of wave-particle duality.
Experiments demonstrated that de Broglie hypothesis is applicable to atoms and macromolecules.
The uncertainty principle is inherent in the properties of all wave-like systems, and that it arises in quantum mechanics is
simply due to the matter wave nature of all quantum objects.
The uncertainty principle is not a statement about the observational success of current technology.
The more precisely the position of some particle is determined, the less precisely its momentum can be known, and vice versa.
¯¯
¯
According to the Copenhagen interpretation, the probabilistic nature of quantum mechanics is intrinsic in our physical universe.
When quantum wave function collapse occurs, physical possibilities are reduced into a single possibility as seen by an observer.
Once a particle in an entangled state is measured and its state is determined, the Copenhagen interpretation demands that the
state of the other entangled particle is also determined instantaneously.
Key Terms
black body radiation: The type of electromagnetic radiation within or surrounding a body in thermodynamic equilibrium with
its environment, or emitted by a black body (an opaque and non-reflective body) held at constant, uniform temperature.
photoelectron: Electrons emitted from matter by absorbing energy from electromagnetic radiation.
wave-particle duality: A postulation that all particles exhibit both wave and particle properties. It is a central concept of
quantum mechanics.
Planck constant: a physical constant that is the quantum of action in quantum mechanics. It has a unit of angular momentum.
The Planck constant was first described as the proportionality constant between the energy of a photon (unit of electromagnetic
radiation) and the frequency of its associated electromagnetic wave in his derivation of the Planck’s law
Maxwell’s equations: A set of equations describing how electric and magnetic fields are generated and altered by each other
and by charges and currents.
elementary particle: a particle not known to have any substructure
photoelectric effect: The occurrence of electrons being emitted from matter (metals and non-metallic solids, liquids, or gases)
as a consequence of their absorption of energy from electromagnetic radiation.
cryptography: the practice and study of techniques for secure communication in the presence of third parties
relativistic quantum mechanics: a theoretical framework for constructing quantum mechanical models of fields and many-
body systems
string theory: an active research framework in particle physics that attempts to reconcile quantum mechanics and general
relativity
black body: An idealized physical body that absorbs all incident electromagnetic radiation, regardless of frequency or angle of
incidence. Although black body is a theoretical concept, you can find approximate realizations of black body in nature.
grating: Any regularly spaced collection of essentially identical, parallel, elongated elements.
Schrödinger equation: A partial-differential that describes how the quantum state of some physical system changes with time.
It was formulated in late 1925 and published in 1926 by the Austrian physicist Erwin Schrödinger
harmonic oscillator: a system that, when displaced from its equilibrium position, experiences a restoring force F proportional
to the displacement x
diffraction: The bending of a wave around the edges of an opening or an obstacle.
special relativity: A theory that (neglecting the effects of gravity) reconciles the principle of relativity with the observation that
the speed of light is constant in all frames of reference.
wave-particle duality: A postulation that all particles exhibit both wave and particle properties. It is a central concept of
quantum mechanics.
matter wave: A concept reflects the wave-particle duality of matter. The theory was proposed by Louis de Broglie.
Rayleigh criterion: The angular resolution of an optical system can be estimated from the diameter of the aperture and the
wavelength of the light, which was first proposed by Lord Rayleigh.
probability density function: Any function whose integral over a set gives the probability that a random variable has a value in
that set.
Bell’s theorem: A no-go theorem famous for drawing an important line in the sand between quantum mechanics (QM) and the
world as we know it classically. In its simplest form, Bell’s theorem states: No physical theory of local hidden variables can
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ever reproduce all of the predictions of quantum mechanics.
epistemological: Of or pertaining to epistemology or theory of knowledge, as a field of study.
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28.2: Applications of Quantum Mechanics
learning objectives
Compare mechanisms of fluorescence and phosphorescence light emission
Fluorescence: Fluorescent minerals emit visible light when exposed to ultraviolet light
Phosphorescence is a specific type of photoluminescence related to fluorescence. Unlike fluorescence, a phosphorescent material
does not immediately re-emit the radiation it absorbs. Excitation of electrons to a higher state is accompanied with the change of a
spin state. Once in a different spin state, electrons cannot relax into the ground state quickly because the re-emission involves
quantum mechanically forbidden energy state transitions. As these transitions occur very slowly in certain materials, absorbed
radiation may be re-emitted at a lower intensity for up to several hours after the original excitation.
Fluorescence and Phosphorescence: Energy scheme used to explain the difference between fluorescence and phosphorescence
Commonly seen examples of phosphorescent materials are the glow-in-the-dark toys, paint, and clock dials that glow for some time
after being charged with a bright light such as in any normal reading or room light. Typically the glowing then slowly fades out
within minutes (or up to a few hours) in a dark room.
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Phosphorescence: Phosphorescent material glowing in the dark.
Lasers
A laser is a device that emits monochromatic light through a process of optical amplification based on the stimulated emission of
photons.
learning objectives
Identify process that generates laser emission and the defining characteristics of laser light
A laser is a device that emits monochromatic light (electromagnetic radiation ). It does so through a process of optical amplification
based on the stimulated emission of photons. The term “laser” originated as an acronym for Light Amplification by Stimulated
Emission of Radiation. Laser is distinct from other light sources for its high degree of spatial and temporal coherence, which means
that laser outputs a narrow beam that maintains its temporal-phase relationship.
Principles of laser operation are largely based on quantum mechanics. (One exception would be free-electron lasers, whose
operation can be explained solely by classical electrodynamics. ) When an electron is excited from a lower-energy to a higher-
energy level, it will not stay that way forever. An electron in an excited state may decay to an unoccupied lower-energy state
according to a particular time constant characterizing that transition. When such an electron decays without external influence, it
emits a photon; this process is called “spontaneous emission. ” The phase associated with the emitted photon is random. A material
with many atoms in an excited state may thus result in radiation that is very monochromatic, but the individual photons would have
no common phase relationship and would emanate in random directions. This is the mechanism of fluorescence and thermal
emission.
However, an external photon at a frequency associated with the atomic transition can affect the quantum mechanical state of the
atom. As the incident photon passes by, the rate of transitions of the excited atom can be significantly enhanced beyond that due to
spontaneous emission. This “induced” decay process is called stimulated emission. In stimulated emission, the decaying atom
produces an identical “copy” of the incoming photon. Therefore, after the atom decays, we have two identical outgoing photons.
Since there was only one incoming photon, we amplified the intensity of light by a factor of 2!
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Stimulated Photon Emission: In stimulated emission process, a photon (with a frequency equal to the atomic transition)
encounters an excited atom, and a new photon identical to the incoming photon is produced. The result is an atom in the ground
state with two outgoing photons.
Holography
Holography is an optical technique which enables three-dimensional images to be made.
learning objectives
Explain how holographic images are recorded and their properties
Holography is a technique which enables three-dimensional images to be made. It involves the use of a laser, interference,
diffraction, light intensity recording and suitable illumination of the recording. The image changes as the position and orientation of
the viewing system changes in exactly the same way as if the object were still present, thus making the image appear three-
dimensional.
Laser: Holograms are recorded using a flash of light that illuminates a scene and then imprints on a recording medium, much in the
way a photograph is recorded. In addition, however, part of the light beam must be shone directly onto the recording medium – this
second light beam is known as the reference beam (]). A hologram requires a laser as the sole light source. Laser is required as a
light source to produce an interference pattern on the recording plate. To prevent external light from interfering, holograms are
usually taken in darkness, or in low level light of a different color from the laser light used in making the hologram. Holography
requires a specific exposure time, which can be controlled using a shutter, or by electronically timing the laser
Recording a hologram: Holograms are recorded using a flash of light that illuminates a scene and then imprints on a recording
medium, much in the way a photograph is recorded. In addition, however, part of the light beam must be shone directly onto the
recording medium – this second light beam is known as the reference beam.
Apparatus: A hologram can be made by shining part of the light beam directly onto the recording medium, and the other part onto
the object in such a way that some of the scattered light falls onto the recording medium. A more flexible arrangement for
recording a hologram requires the laser beam to be aimed through a series of elements that change it in different ways. The first
element is a beam splitter that divides the beam into two identical beams, each aimed in different directions:
One beam (known as the illumination or object beam) is spread using lenses and directed onto the scene using mirrors. Some of
the light scattered (reflected) from the scene then falls onto the recording medium.
The second beam (known as the reference beam) is also spread through the use of lenses, but is directed so that it doesn’t come
in contact with the scene, and instead travels directly onto the recording medium.
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Reconstructing a hologram: An interference pattern can be considered an encoded version of a scene, requiring a particular key –
the original light source – in order to view its contents. This missing key is provided later by shining a laser, identical to the one
used to record the hologram, onto the developed film. When this beam illuminates the hologram, it is diffracted by the hologram’s
surface pattern. This produces a light field identical to the one originally produced by the scene and scattered onto the hologram
Several different materials can be used as the recording medium. One of the most common is a film very similar to photographic
film (silver halide photographic emulsion), but with a much higher concentration of light-reactive grains, making it capable of the
much higher resolution that holograms require. A layer of this recording medium (e.g. silver halide) is attached to a transparent
substrate, which is commonly glass, but may also be plastic.
Process: When the two laser beams reach the recording medium, their light waves intersect and interfere with each other. It is this
interference pattern that is imprinted on the recording medium. The pattern itself is seemingly random, as it represents the way in
which the scene’s light interfered with the original light source – but not the original light source itself. The interference pattern can
be considered an encoded version of the scene, requiring a particular key – the original light source – in order to view its contents.
This missing key is provided later by shining a laser, identical to the one used to record the hologram, onto the developed film.
When this beam illuminates the hologram, it is diffracted by the hologram’s surface pattern. This produces a light field identical to
the one originally produced by the scene and scattered onto the hologram. The image this effect produces in a person’s retina is
known as a virtual image.
learning objectives
Explain how properties of elements vary within groups and across periods in the periodic table
The periodic table is a tabular display of the chemical elements. The elements are organized based on their atomic numbers,
electron configurations, and recurring chemical properties.
In the periodic table, elements are presented in order of increasing atomic number (the number of protons). The rows of the table
are called periods; the columns of the s- (columns 1-2 and He), d- (columns 3-12), and p-blocks (columns 13-18, except He) are
called groups. (The terminology of s-, p-, and d- blocks originate from the valence atomic orbitals the element’s electrons occupy. )
Some groups have specific names, such as the halogens or the noble gases. Since, by definition, a periodic table incorporates
recurring trends, any such table can be used to derive relationships between the properties of the elements and predict the properties
of new, yet-to-be-discovered, or synthesized elements. As a result, the periodic table provides a useful framework for analyzing
chemical behavior, and such tables are widely used in chemistry and other sciences.
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Blocks in the Periodic Table: A diagram of the periodic table, highlighting the different blocks
Groups
Agroup, or family, is a vertical column in the periodic table. Groups usually have more significant periodic trends than do periods
and blocks, which are explained below. Modern quantum mechanical theories of atomic structure explain group trends by
proposing that elements in the same group generally have the same electron configurations in their valence (or outermost, partially
filled) shell. Consequently, elements in the same group tend to have shared chemistry and exhibit a clear trend in properties with
increasing atomic number. However, in some parts of the periodic table, such as the d-block and the f-block, horizontal similarities
can be as important as, or more pronounced than, vertical similarities.
Periods
A period is a horizontal row in the periodic table. Although groups generally have more significant periodic trends, there are
regions where horizontal trends are more significant than vertical group trends, such as in the f-block, where the lanthanides and
actinides form two substantial horizontal series of elements. Elements in the same period show trends in atomic radius, ionization
energy, and electron affinity. Atomic radius usually decreases from left to right across a period. This occurs because each
successive element has an added proton and electron, which causes the electron to be drawn closer to the nucleus, decreasing the
radius.
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The periodic table: Here is the complete periodic table with atomic numbers, groups, and periods. Each entry on the periodic table
represents one element, and compounds are made up of several of these elements.
X-Rays
X-rays are a form of electromagnetic radiation and have wavelengths in the range of 0.01 to 10 nanometers.
learning objectives
Describe the properties of X-rays and how can be generated
X-radiation (composed of x-rays) is a form of electromagnetic radiation. X-rays have wavelengths in the range of 0.01 to 10
nanometers, which corresponds to frequencies in the range of 30 petahertz to 30 exahertz (3·1016 Hz to 3·1019 Hz) and energies in
the of range 100 eV to 100 keV.
X-Ray Spectrum and Applications: X-rays are part of the electromagnetic spectrum, with wavelengths shorter than those of
visible light. Different applications use different parts of the X-ray spectrum.
X-rays can be generated by an x-ray tube, a vacuum tube that uses high voltage to accelerate the electrons released by a hot cathode
to a high velocity. The high-velocity electrons collide with a metal target, the anode, creating the x-rays. The maximum energy of
the produced x-ray photon is limited by the energy of the incident electron, which is equal to the voltage on the tube times the
electron charge, so an 80-kV tube cannot create x-rays with an energy greater than 80 keV. When the electrons hit the target, x-rays
are created through two different atomic processes:
1. X-ray fluorescence, if the electron has enough energy that it can knock an orbital electron out of the inner electron shell of a
metal atom. As a result, electrons from higher energy levels fill up the vacancy, and x-ray photons are emitted. This process
produces an emission spectrum of x-rays at a few discrete frequencies, sometimes referred to as the spectral lines. The spectral
lines generated depend on the target (anode) element used and therefore are called characteristic lines. Usually these are
transitions from upper shells into the K shell (called K lines), or the L shell (called L lines), and so on.
2. Bremsstrahlung, literally meaning braking radiation. Bremsstrahlung is radiation given off by the electrons as they are scattered
by the strong electric field near the high-Z (proton number) nuclei. These x-rays have a continuous spectrum. The intensity of
the x-rays increases linearly with decreasing frequency, from zero at the energy of the incident electrons, the voltage on the x-
ray tube.
Both of these x-ray production processes are inefficient, with a production efficiency of only about one percent. Therefore, to
produce a usable flux of x-rays, most of the electric power consumed by the tube is released as heat waste. The x-ray tube must be
designed to dissipate this excess heat.
A specialized source of x-rays that is becoming widely used in research is synchrotron radiation, which is generated by particle
accelerators. Its unique features are x-ray outputs many orders of magnitude greater than those of x-ray tubes, wide x-ray spectra,
excellent collimation, and linear polarization.
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learning objectives
Identify major contributions to the understanding of atomic structure that were made by Niels Bohr, Erwin Schrödinger,
and Werner Heisenberg
The atom is a basic unit of matter that consists of a nucleus surrounded by negatively charged electrons. The atomic nucleus
contains a mix of positively charged protons and electrically neutral neutrons. The electrons of an atom are bound to the nucleus by
the electromagnetic (Coulomb) force. Atoms are minuscule objects with diameters of a few tenths of a nanometer and tiny masses
proportional to the volume implied by these dimensions. Atoms in solid states (or, to be precise, their electron clouds) can be
observed individually using special instruments such as the scanning tunneling microscope.
Hydrogen-1 (one proton + one electron) is the simplest form of atoms, and not surprisingly, our quantum mechanical understanding
of atoms evolved with the understanding of this species. In 1913, physicist Niels Bohr suggested that the electrons were confined
into clearly defined, quantized orbits, and could jump between these, but could not freely spiral inward or outward in intermediate
states. An electron must absorb or emit specific amounts of energy to transition between these fixed orbits. Bohr’s model
successfully explained spectroscopic data of hydrogen very well, but it adopted a semiclassical approach where electron was still
considered a (classical) particle.
Adopting Louis de Broglie’s proposal of wave-particle duality, Erwin Schrödinger, in 1926, developed a mathematical model of the
atom that described the electrons as three-dimensional waveforms rather than point particles. A consequence of using waveforms to
describe particles is that it is mathematically impossible to obtain precise values for both the position and momentum of a particle
at the same time; this became known as the uncertainty principle, formulated by Werner Heisenberg in 1926. Thereafter, the
planetary model of the atom was discarded in favor of one that described atomic orbital zones around the nucleus where a given
electron is most likely to be observed.
Illustration of the Helium Atom: This is an illustration of the helium atom, depicting the nucleus (pink) and the electron cloud
distribution (black). The nucleus (upper right) in helium-4 is in reality spherically symmetric and closely resembles the electron
cloud, although for more complicated nuclei this is not always the case. The black bar is one angstrom (10-10 m, or 100 pm).
Modern quantum mechanical view of hydrogen has evolved further after Schrödinger, by taking relativistic correction terms into
account. Quantum electrodynamics (QED), a relativistic quantum field theory describing the interaction of electrically charged
particles, has successfully predicted minuscule corrections in energy levels. One of the hydrogen’s atomic transitions (n=2 to n=1,
n: principal quantum number) has been measured to an extraordinary precision of 1 part in a hundred trillion. This kind of
spectroscopic precision allows physicists to refine quantum theories of atoms, by accounting for minuscule discrepancies between
experimental results and theories.
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Key Points
The emitted light usually has a longer wavelength, and therefore lower energy, than the absorbed radiation.
Fluorescence occurs when an orbital electron of a molecule or atom relaxes to its ground state by emitting a photon of light
after being excited to a higher quantum state by some type of energy.
In a phosphorescence, excitation of electrons to a higher state is accompanied with the change of a spin state. Relaxation is a
slow process since it involves energy state transitions “forbidden” in quantum mechanics.
Principles of laser operation are largely based on quantum mechanics, most importantly on the process of the stimulated
emission of photons.
Spontaneous emission is a random decaying process. The phase associated with the emitted photon is also random.
Atomic transition can be stimulated by the presence of an incoming photon at a frequency associated with the atomic transition.
This process leads to optical amplification as an identical photon is emitted along with the incoming photon.
When the two laser beams reach the recording medium, their light waves intersect and interfere with each other. It is this
interference pattern that is imprinted on the recording medium.
When a reconstruction beam illuminates the hologram, it is diffracted by the hologram’s surface pattern. This produces a light
field identical to the one originally produced by the scene and scattered onto the hologram.
Holographic image changes as the position and orientation of the viewing system changes in exactly the same way as if the
object were still present, thus making the image appear three-dimensional.
A periodic table is a useful framework for analyzing chemical behavior. Such tables are widely used in chemistry and other
sciences.
A group, or family, is a vertical column in the periodic table. Groups usually have more significant periodic trends than do
periods and blocks.
A period is a horizontal row in the periodic table. Elements in the same period show trends in atomic radius, ionization energy,
electron affinity, and electronegativity.
X-rays can be generated by an x-ray tube, a vacuum tube, or a particle accelerator.
X-ray fluorescence and Bremsstrahlung are processes through which x-rays are produced.
Synchrotron radiation is generated by particle accelerators. Its unique features are x-ray outputs many orders of magnitude
greater than those of x-ray tubes, wide x-ray spectra, excellent collimation, and linear polarization.
Niels Bohr suggested that the electrons were confined into clearly defined, quantized orbits, and could jump between these, but
could not freely spiral inward or outward in intermediate states.
Erwin Schrödinger, in 1926, developed a mathematical model of the atom that described the electrons as three-dimensional
waveforms rather than point particles.
Modern quantum mechanical view of hydrogen has evolved further after Schrödinger, by taking relativistic correction terms
into account. This is referred to a quantum electrodynamics (QED).
Key Terms
spin: A quantum angular momentum associated with subatomic particles; it also creates a magnetic moment.
photon: The quantum of light and other electromagnetic energy, regarded as a discrete particle having zero rest mass, no
electric charge, and an indefinitely long lifetime.
ground state: the stationary state of lowest energy of a particle or system of particles
free-electron laser: a laser that use a relativistic electron beam as the lasing medium, which moves freely through a magnetic
structure
monochromatic: Describes a beam of light with a single wavelength (i.e., of one specific color or frequency).
coherence: an ideal property of waves that enables stationary (i.e., temporally and spatially constant) interference
interference: An effect caused by the superposition of two systems of waves, such as a distortion on a broadcast signal due to
atmospheric or other effects.
laser: A device that produces a monochromatic, coherent beam of light.
silver halide: The light-sensitive chemicals used in photographic film and pape
atomic orbital: The quantum mechanical behavior of an electron in an atom describing the probability of the electron’s
particular position and energy.
electron affinity: the amount of energy released when an electron is added to a neutral atom or molecule to form a negative ion
ionization energy: the amount of energy required to remove an electron from an atom or molecule in the gas phase
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photon: The quantum of light and other electromagnetic energy, regarded as a discrete particle having zero rest mass, no
electric charge, and an indefinitely long lifetime.
particle accelerator: A device that accelerates electrically charged particles to extremely high speeds, for the purpose of
inducing high-energy reactions or producing high-energy radiation.
wave-particle duality: A postulation that all particles exhibit both wave and particle properties. It is a central concept of
quantum mechanics.
scanning tunneling microscope: An instrument for imaging surfaces at the atomic level.
semiclassical approach: A theory in which one part of a system is described quantum-mechanically whereas the other is
treated classically.
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CHAPTER OVERVIEW
29: Atomic Physics is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
29.1: Overview
learning objectives
Identify assumption made by Max Planck to describe the electromagnetic radiation emitted by a black body
A black body in thermal equilibrium (i.e. at a constant temperature) emits electromagnetic radiation called black body radiation.
Black body radiation has a characteristic, continuous frequency spectrum that depends only on the body’s temperature. Max
Planck, in 1901, accurately described the radiation by assuming that electromagnetic radiation was emitted in discrete packets (or
quanta). Planck’s quantum hypothesis is a pioneering work, heralding advent of a new era of modern physics and quantum theory.
Explaining the properties of black-body radiation was a major challenge in theoretical physics during the late nineteenth century.
Predictions based on classical theories failed to explain black body spectra observed experimentally, especially at shorter
wavelength. The puzzle was solved in 1901 by Max Planck in the formalism now known as Planck’s law of black-body radiation.
Contrary to the common belief that electromagnetic radiation can take continuous values of energy, Planck introduced a radical
concept that electromagnetic radiation was emitted in discrete packets (or quanta) of energy. Although Planck’s derivation is
beyond the scope of this section (it will be covered in Quantum Mechanics), Planck’s law may be written:
2
2hc 1
Bλ (T) = (29.1.1)
bc
5
λ λk T
e B −1
where BB is the spectral radiance of the surface of the black body, TT is its absolute temperature, λλ is wavelength of the radiation,
kBkB is the Boltzmann constant, hh is the Planck constant, and cc is the speed of light. This equation explains the black body
spectra shown below. Planck’s quantum hypothesis is one of the breakthroughs in the modern physics. It is not a surprise that he
introduced Planck constant h = 6.626 × 10 J ⋅ s for the first time in his derivation of the Planck’s law.
−34
UV VISIBLE INFRARED
14
5000 K
Spectral radiance (kW · sr⁻¹ · m⁻² · nm⁻¹)
12
Classical theory (5000 K)
10
4000 K
4
2
3000 K
0
0 0.5 1 1.5 2 2.5 3
Wavelength (μm)
Black body radiation spectrum: Typical spectrum from a black body at different temperatures (shown in blue, green and red
curves). As the temperature decreases, the peak of the black-body radiation curve moves to lower intensities and longer
wavelengths. Black line is a prediction of a classical theory for an object at 5,000K, showing catastropic discrepancy at shorter
wavelengh.
Note that the spectral radiance depends on two variables, wavelength and temperature. The radiation has a specific spectrum and
intensity that depends only on the temperature of the body. Despite its simplicity, Planck’s law describes radiation properties of
objects (e.g. our body, planets, stars) reasonably well.
Key Points
A black body in thermal equilibrium emits electromagnetic radiation called black body radiation.
The radiation has a specific spectrum and intensity that depends only on the temperature of the body.
Max Planck, in 1901, accurately described the radiation by assuming that electromagnetic radiation was emitted in discrete
packets (or quanta). Planck’s quantum hypothesis is a pioneering work, heralding advent of a new era of modern physics and
quantum theory.
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Key Terms
spectral radiance: measures of the quantity of radiation that passes through or is emitted from a surface and falls within a
given solid angle in a specified direction.
black body: An idealized physical body that absorbs all incident electromagnetic radiation, regardless of frequency or angle of
incidence. Although black body is a theoretical concept, you can find approximate realizations of black body in nature.
Planck constant: a physical constant that is the quantum of action in quantum mechanics. It has a unit of angular momentum.
The Planck constant was first described as the proportionality constant between the energy of a photon (unit of electromagnetic
radiation) and the frequency of its associated electromagnetic wave in his derivation of the Planck’s law
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29.2: The Early Atom
learning objectives
Discuss experiments that led to discovery of the electron and the nucleus
Though originally viewed as a particle that cannot be cut into smaller particles, modern scientific usage denotes the atom as
composed of various subatomic particles. The constituent particles of an atom (each discovered independently) are: the electron,
the proton and the neutron. (The hydrogen-1 atom, however, has no neutrons, and a positive hydrogen ion has no electrons. )
Classical Atomic Model: Atomic model before the advent of Quantum Mechanics.
Electron
The German physicist Johann Wilhelm Hittorf undertook the study of electrical conductivity in rarefied gases. In 1869, he
discovered a glow emitted from the cathode that increased in size with decrease in gas pressure. In 1896, the British physicist J. J.
Thomson performed experiments demonstrating that cathode rays were unique particles, rather than waves, atoms or molecules, as
was believed earlier. Thomson made good estimates of both the charge ee and the mass mm, finding that cathode ray particles
(which he called “corpuscles”) had perhaps one thousandth the mass of hydrogen, the least massive ion known. He showed that
their charge to mass ratio (e/m) was independent of cathode material. (Fig 1 shows a beam of deflected electrons. )
Proton
In 1917 (in experiments reported in 1919), Rutherford proved that the hydrogen nucleus is present in other nuclei, a result usually
described as the discovery of the proton. Earlier, Rutherford learned to create hydrogen nuclei as a type of radiation produced as a
yield of the impact of alpha particles on hydrogen gas; these nuclei were recognized by their unique penetration signature in air and
their appearance in scintillation detectors. These experiments began when Rutherford noticed that when alpha particles were shot
into air (mostly nitrogen), his scintillation detectors displayed the signatures of typical hydrogen nuclei as a product. After
experimentation Rutherford traced the reaction to the nitrogen in air, and found that the effect was larger when alphas were
produced into pure nitrogen gas. Rutherford determined that the only possible source of this hydrogen was the nitrogen, and
therefore nitrogen must contain hydrogen nuclei. One hydrogen nucleus was knocked off by the impact of the alpha particle,
producing oxygen-17 in the process. This was the first reported nuclear reaction, 14N + α → 17O + p .
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Neutron
In 1920, Ernest Rutherford conceived the possible existence of the neutron. In particular, Rutherford examined the disparity found
between the atomic number of an atom and its atomic mass. His explanation for this was the existence of a neutrally charged
particle within the atomic nucleus. He considered the neutron to be a neutral double consisting of an electron orbiting a proton. In
1932, James Chadwick showed uncharged particles in the radiation he used. These particles had a similar mass as protons, but did
not have the same characteristics as protons. Chadwick followed some of the predictions of Rutherford, the first to work in this
then unknown field.
learning objectives
Describe postulates of Dalton’s atomic theory and the atomic theories of ancient Greek philosophers
The atom is a basic unit of matter that consists of a dense central nucleus surrounded by a cloud of negatively charged electrons.
The atomic nucleus contains a mix of positively charged protons and electrically neutral neutrons (except in the case of hydrogen-
1, which is the only stable nuclide with no neutrons). The electrons of an atom are bound to the nucleus by the electromagnetic
force. We have a detailed (and accurate) model of the atom now, but it took a long time to come up with the correct answer.
Illustration of the Helium Atom: This is an illustration of the helium atom, depicting the nucleus (pink) and the electron cloud
distribution (black). The nucleus (upper right) in helium-4 is in reality spherically symmetric and closely resembles the electron
cloud, although for more complicated nuclei this is not always the case. The black bar is one angstrom (10−1010−10 m, or 100
pm).
People have long speculated about the structure of matter and the existence of atoms. The earliest significant ideas to survive are
from the ancient Greeks in the fifth century BC, especially from the philosophers Leucippus and Democritus. (There is some
evidence that philosophers in both India and China made similar speculations at about the same time. ) They considered the
question of whether a substance can be divided without limit into ever smaller pieces. There are only a few possible answers to this
question. One is that infinitesimally small subdivision is possible. Another is what Democritus in particular believed — that there is
a smallest unit that cannot be further subdivided. Democritus called this the atom. We now know that atoms themselves can be
subdivided, but their identity is destroyed in the process, so the Greeks were correct in a respect. The Greeks also felt that atoms
were in constant motion, another correct notion.
The Greeks and others speculated about the properties of atoms, proposing that only a few types existed and that all matter was
formed as various combinations of these types. The famous proposal that the basic elements were earth, air, fire, and water was
brilliant but incorrect. The Greeks had identified the most common examples of the four states of matter (solid, gas, plasma, and
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liquid) rather than the basic chemical elements. More than 2000 years passed before observations could be made with equipment
capable of revealing the true nature of atoms.
Over the centuries, discoveries were made regarding the properties of substances and their chemical reactions. Certain systematic
features were recognized, but similarities between common and rare elements resulted in efforts to transmute them (lead into gold,
in particular) for financial gain. Secrecy was commonplace. Alchemists discovered and rediscovered many facts but did not make
them broadly available. As the Middle Ages ended, the practice of alchemy gradually faded, and the science of chemistry arose. It
was no longer possible, nor considered desirable, to keep discoveries secret. Collective knowledge grew, and by the beginning of
the 19th century, an important fact was well established: the masses of reactants in specific chemical reactions always have a
particular mass ratio. This is very strong indirect evidence that there are basic units (atoms and molecules) that have these same
mass ratios. English chemist John Dalton (1766-1844) did much of this work, with significant contributions by the Italian physicist
Amedeo Avogadro (1776-1856). It was Avogadro who developed the idea of a fixed number of atoms and molecules in a mole.
This special number is called Avogadro’s number in his honor (6.022 ⋅ 1023).
Dalton believed that matter is composed of discrete units called atoms, as opposed to the obsolete notion that matter could be
divided into any arbitrarily small quantity. He also believed that atoms are the indivisible, ultimate particles of matter. However,
this belief was overturned near the end of the 19th century by Thomson, with his discovery of electrons.
Intro to the History of Atomic Theory – Intro: Rutherford, Thomson, electrons, nuclei, and plums. I don’t mean to be a bohr, but
do you think pudding should have a role in serious scientific inquiry?
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learning objectives
Describe model of an atom proposed by J. J. Thomson.
J. J. Thomson, who discovered the electron in 1897, proposed the plum pudding model of the atom in 1904 before the discovery of
the atomic nucleus in order to include the electron in the atomic model. In Thomson’s model, the atom is composed of electrons
(which Thomson still called “corpuscles,” though G. J. Stoney had proposed that atoms of electricity be called electrons in 1894)
surrounded by a soup of positive charge to balance the electrons’ negative charges, like negatively charged “plums” surrounded by
positively charged “pudding”. The electrons (as we know them today) were thought to be positioned throughout the atom in
rotating rings. In this model the atom was also sometimes described to have a “cloud” of positive charge.
Plum pudding model of the atom: A schematic presentation of the plum pudding model of the atom; in Thomson’s mathematical
model the “corpuscles” (in modern language, electrons) were arranged non-randomly, in rotating rings.
With this model, Thomson abandoned his earlier “nebular atom” hypothesis, in which the atom was composed of immaterial
vortices. Now, at least part of the atom was to be composed of Thomson’s particulate negative corpuscles, although the rest of the
positively charged part of the atom remained somewhat nebulous and ill-defined.
The 1904 Thomson model was disproved by the 1909 gold foil experiment performed by Hans Geiger and Ernest Marsden. This
gold foil experiment was interpreted by Ernest Rutherford in 1911 to suggest that there is a very small nucleus of the atom that
contains a very high positive charge (in the case of gold, enough to balance the collective negative charge of about 100 electrons).
His conclusions led him to propose the Rutherford model of the atom.
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Intro to the History of Atomic Theory – The Thomson Model: Rutherford, Thomson, electrons, nuclei, and plums. I don’t mean
to be a bohr, but do you think pudding should have a role in serious scientific inquiry?
learning objectives
Describe gold foil experiment performed by Geiger and Marsden under directions of Rutherford and its implications for the
model of the atom
The Rutherford model is a model of the atom named after Ernest Rutherford. Rutherford directed the famous Geiger-Marsden
experiment in 1909, which suggested, according to Rutherford’s 1911 analysis, that J. J. Thomson’s so-called “plum pudding
model” of the atom was incorrect. Rutherford’s new model for the atom, based on the experimental results, contained the new
features of a relatively high central charge concentrated into a very small volume in comparison to the rest of the atom. This central
volume also contained the bulk of the atom’s mass. This region would later be named the “nucleus. ”
Atomic Planetary Model: Basic diagram of the atomic planetary model; electrons are in green, and the nucleus is in red
In 1911, Rutherford designed an experiment to further explore atomic structure using the alpha particles emitted by a radioactive
element. Following his direction, Geiger and Marsden shot alpha particles with large kinetic energies toward a thin foil of gold.
Measuring the pattern of scattered particles was expected to provide information about the distribution of charge within the atom.
Under the prevailing plum pudding model, the alpha particles should all have been deflected by, at most, a few degrees. However,
the actual results surprised Rutherford. Although many of the alpha particles did pass through as expected, many others were
deflected at small angles while others were reflected back to the alpha source.
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From purely energetic considerations of how far particles of known speed would be able to penetrate toward a central charge of 100
e, Rutherford was able to calculate that the radius of his gold central charge would need to be less than 3.4⋅10−143.4⋅10−14 meters.
This was in a gold atom known to be about 10−1010−10 meters in radius; a very surprising finding, as it implied a strong central
charge less than 1300013000th of the diameter of the atom.
Intro to the History of Atomic Theory – The Rutherford Model: Rutherford, Thomson, electrons, nuclei, and plums. I don’t
mean to be a bohr, but do you think pudding should have a role in serious scientific inquiry?
learning objectives
Describe model of atom proposed by Niels Bohr.
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Niels Bohr: Niels Bohr, Danish physicist, used the planetary model of the atom to explain the atomic spectrum and size of the
hydrogen atom. His many contributions to the development of atomic physics and quantum mechanics; his personal influence on
many students and colleagues; and his personal integrity, especially in the face of Nazi oppression, earned him a prominent place in
history. (credit: Unknown Author, via Wikimedia Commons)
For decades, many questions had been asked about atomic characteristics. From their sizes to their spectra, much was known about
atoms, but little had been explained in terms of the laws of physics. Bohr’s theory explained the atomic spectrum of hydrogen,
made him instantly famous, and established new and broadly applicable principles in quantum mechanics.
One big puzzle that the planetary-model of atom had was the following. The laws of classical mechanics predict that the electron
should release electromagnetic radiation while orbiting a nucleus (according to Maxwell’s equations, accelerating charge should
emit electromagnetic radiation). Because the electron would lose energy, it would gradually spiral inwards, collapsing into the
nucleus. This atom model is disastrous, because it predicts that all atoms are unstable. Also, as the electron spirals inward, the
emission would gradually increase in frequency as the orbit got smaller and faster. This would produce a continuous smear, in
frequency, of electromagnetic radiation. However, late 19th century experiments with electric discharges have shown that atoms
will only emit light (that is, electromagnetic radiation) at certain discrete frequencies.
To overcome this difficulty, Niels Bohr proposed, in 1913, what is now called the Bohr model of the atom. He suggested that
electrons could only have certain classical motions:
1. Electrons in atoms orbit the nucleus.
2. The electrons can only orbit stably, without radiating, in certain orbits (called by Bohr the “stationary orbits”): at a certain
discrete set of distances from the nucleus. These orbits are associated with definite energies and are also called energy shells or
energy levels. In these orbits, the electron’s acceleration does not result in radiation and energy loss as required by classical
electrodynamics.
3. Electrons can only gain and lose energy by jumping from one allowed orbit to another, absorbing or emitting electromagnetic
radiation with a frequency νν determined by the energy difference of the levels according to the Planck relation:
ΔE = E2 − E1 = hν (29.2.1)
Semiclassical Model
The significance of the Bohr model is that the laws of classical mechanics apply to the motion of the electron about the nucleus
only when restricted by a quantum rule. Therefore, his atomic model is called a semiclassical model.
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learning objectives
Describe basic assumptions that were applied by Niels Bohr to the planetary model of an atom
In previous modules, we have seen puzzles from classical atomic theories (e.g., the Rutherford model). Most importantly, classical
electrodynamics predicts that an atom described by a (classical) planetary model would be unstable. To explain the puzzle, Bohr
proposed what is now called the Bohr model of the atom in 1913. He suggested that electrons could only have certain classical
motions:
1. Electrons in atoms orbit the nucleus.
2. The electrons can only orbit stably, without radiating, in certain orbits (called by Bohr the “stationary orbits”) at a certain
discrete set of distances from the nucleus. These orbits are associated with definite energies and are also called energy shells or
energy levels. In these orbits, the electron’s acceleration does not result in radiation and energy loss as required by classical
electrodynamics.
3. Electrons can only gain and lose energy by jumping from one allowed orbit to another, absorbing or emitting electromagnetic
radiation with a frequency νν determined by the energy difference of the levels according to the Planck relation:
ΔE = E − E = hν , where h is the Planck constant. In addition, Bohr also assumed that the angular momentum L is
2 1
restricted to be an integer multiple of a fixed unit: L = n = nℏ , where n = 1, 2, 3, … is called the principal quantum
h
2π
number, and ℏ = h
2π
.
We have seen that Planck adopted a new condition of energy quantization to explain the black body radiation, where he introduced
the Planck constant hh for the first time. Soon after, Einstein resorted to this new concept of energy quantization and used the
Planck constant again to explain the photoelectric effects, in which he assumed that electromagnetic radiation interact with matter
as particles (later named “photons”). Here, Bohr explained the atomic hydrogen spectrum successfully for the first time by adopting
a quantization condition and by introducing the Planck constant in his atomic model. Over the period of radical development in the
early 20th century, physicists began to realize that it was essential to introduce the notion of “quantization” to explain microscopic
worlds.
Rutherford-Bohr model: The Rutherford–Bohr model of the hydrogen atom (Z = 1 ) or a hydrogen-like ion (Z > 1 ), where the
negatively charged electron confined to an atomic shell encircles a small, positively charged atomic nucleus, and where an electron
jump between orbits is accompanied by an emitted or absorbed amount of electromagnetic energy (hνhν). The orbits in which the
electron may travel are shown as gray circles; their radius increases as n , where n is the principal quantum number. The 3 → 2
2
transition depicted here produces the first line of the Balmer series, and for hydrogen (Z = 1 ) it results in a photon of wavelength
656 nm (red light).
Bohr Orbits
According to Bohr, electrons can only orbit stably, in certain orbits, at a certain discrete set of distances from the nucleus.
learning objectives
Explain relationship between the “Bohr orbits” and the quantization effect
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Danish Physicist Neils Bohr was clever enough to discover a method of calculating the electron orbital energies in hydrogen. As
we’ve seen in the previous module “The Bohr Model of Atom,” Bohr assumed that the electrons can only orbit stably, without
radiating, in certain orbits (named by Bohr as “stationary orbits”), at a certain discrete set of distances from the nucleus. These
“Bohr orbits” have a very important feature of quantization as shown in the following. This was an important first step that has
been improved upon, but it is well worth repeating here, as it correctly describes many characteristics of hydrogen. Assuming
circular orbits, Bohr proposed that the angular momentum LL of an electron in its orbit is quantized, that is, has only specific,
discrete values. The value for LL is given by the formula:
h
L = me vrn = n = nℏ (29.2.2)
2π
where LL is the angular momentum, meme is the electron’s mass, rnrn is the radius of the nn-th orbit, and hh is Planck’s constant.
Note that angular momentum is L = Iω . For a small object at a radius r, I = mr and ω = , so that:
2 v
v
2
L = (mr ) ( ) = mvr (29.2.3)
r
Quantization says that this value of mvrmvr can only have discrete values. At the time, Bohr himself did not know why angular
momentum should be quantized, but using this assumption he was able to calculate the energies in the hydrogen spectrum,
something no one else had done at the time.
Below is an energy-level diagram, which is a convenient way to display energy states—the allowed energy levels of the electron
(as relative to our discussion). Energy is plotted vertically with the lowest or ground state at the bottom and with excited states
above. Given the energies of the lines in an atomic spectrum, it is possible (although sometimes very difficult) to determine the
energy levels of an atom. Energy-level diagrams are used for many systems, including molecules and nuclei. A theory of the atom
or any other system must predict its energies based on the physics of the system.
Energy-Level Diagram Plot: An energy-level diagram plots energy vertically and is useful in visualizing the energy states of a
system and the transitions between them. This diagram is for the hydrogen-atom electrons, showing a transition between two orbits
having energies E and E .
4 2
learning objectives
Apply proper equation to calculate energy levels and the energy of an emitted photon for a hydrogen-like atom
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From Bohr’s assumptions, we will now derive a number of important properties of the hydrogen atom from the classical physics.
We start by noting the centripetal force causing the electron to follow a circular path is supplied by the Coulomb force. To be more
general, we note that this analysis is valid for any single-electron atom. So, if a nucleus has ZZ protons (Z=1Z=1 for hydrogen,
Z=2Z=2 for helium, etc.) and only one electron, that atom is called a hydrogen-like atom.
The spectra of hydrogen-like ions are similar to hydrogen, but shifted to higher energy by the greater attractive force between the
electron and nucleus. The magnitude of the centripetal force is m2evrnmev2rn, while the Coulomb force is Zkee2r2Zkee2r2. The
tacit assumption here is that the nucleus is more massive than the stationary electron, and the electron orbits about it. This is
consistent with the planetary model of the atom. Equating these:
2 2
me v Zke e
= (29.2.4)
2
r r
The total energy is negative and inversely proportional to rr. This means that it takes energy to pull the orbiting electron away from
the proton. For infinite values of rr, the energy is zero, corresponding to a motionless electron infinitely far from the proton.
Now, here comes the Quantum rule: As we saw in the previous module, the angular momentum L=mervL=merv is an integer
multiple of ℏℏ:
me vr = nℏ (29.2.7)
Substituting the expression in the equation for speed above gives an equation for rr in terms of nn:
−−−−−−−−
2
√Zke e me r = nℏ (29.2.8)
The smallest possible value of rr in the hydrogen atom is called the Bohr radius and is equal to 0.053 nm. The energy of the nn-th
level for any atom is determined by the radius and quantum number:
2 2
2 2 2
Zke e Z (ke e ) me −13.6Z
E =− =− ≈ eV (29.2.10)
2 2 2
2rn 2ℏ n n
Using this equation, the energy of a photon emitted by a hydrogen atom is given by the difference of two hydrogen energy levels:
1 1
E = Ei − Ef = R ( − ) (29.2.11)
2 2
n n
f i
Which is the Rydberg formula describing all the hydrogen spectrum and RR is the Rydberg constant. Bohr’s model predicted
experimental hydrogen spectrum extremely well.
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Fig 1: A schematic of the hydrogen spectrum shows several series named for those who contributed most to their determination.
Part of the Balmer series is in the visible spectrum, while the Lyman series is entirely in the UV, and the Paschen series and others
are in the IR. Values of nf and ni are shown for some of the lines.
Hydrogen Spectra
The observed hydrogen-spectrum wavelengths can be calculated using the following formula: 1
λ
=R(
1
2
n
−
1
2
n
) .
f i
learning objectives
Explain difference between Lyman, Balmer, and Paschen series
For decades, many questions had been asked about atomic characteristics. From their sizes to their spectra, much was known about
atoms, but little had been explained in terms of the laws of physics. Atomic and molecular emission and absorption spectra have
been known for over a century to be discrete (or quantized). Maxwell and others had realized that there must be a connection
between the spectrum of an atom and its structure, something like the resonant frequencies of musical instruments. But, despite
years of efforts by many great minds, no one had a workable theory. (It was a running joke that any theory of atomic and molecular
spectra could be destroyed by throwing a book of data at it, so complex were the spectra.) Following Einstein’s proposal of photons
with quantized energies directly proportional to their wavelengths, it became even more evident that electrons in atoms can exist
only in discrete orbits.
In some cases, it had been possible to devise formulas that described the emission spectra. As you might expect, the simplest atom
—hydrogen, with its single electron—has a relatively simple spectrum. The hydrogen spectrum had been observed in the infrared
(IR), visible, and ultraviolet (UV), and several series of spectral lines had been observed. The observed hydrogen-spectrum
wavelengths can be calculated using the following formula:
1 1 1
= R( − ) (29.2.12)
2 2
λ n n
f i
where λ is the wavelength of the emitted EM radiation and R is the Rydberg constant, determined by the experiment to be
R = 1.097 ⋅ 107 m
−1
, and n , n are positive integers associated with a specific series.
f i
These series are named after early researchers who studied them in particular depth. For the Lyman series, n = 1 for the Balmer
f
series, n = 2 ; for the Paschen series, n = 3 ; and so on. The Lyman series is entirely in the UV, while part of the Balmer series is
f f
visible with the remainder UV. The Paschen series and all the rest are entirely IR. There are apparently an unlimited number of
series, although they lie progressively farther into the infrared and become difficult to observe as nfnfincreases. The constant nini is
a positive integer, but it must be greater than n . Thus, for the Balmer series, n = 2 andn = 3, 4, 5, 6 …. Note that nini can
f f i
approach infinity.
Electron transitions and their resulting wavelengths for hydrogen.: Energy levels are not to scale.
While the formula in the wavelengths equation was just a recipe designed to fit data and was not based on physical principles, it did
imply a deeper meaning. Balmer first devised the formula for his series alone, and it was later found to describe all the other series
by using different values of nfnf. Bohr was the first to comprehend the deeper meaning. Again, we see the interplay between
experiment and theory in physics. Experimentally, the spectra were well established, an equation was found to fit the experimental
data, but the theoretical foundation was missing.
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de Broglie and the Bohr Model
By assuming that the electron is described by a wave and a whole number of wavelengths must fit, we derive Bohr’s quantization
assumption.
learning objectives
Describe reinterpretation of Bohr’s condition by de Broglie
Bohr’s condition, that the angular momentum is an integer multiple of ℏℏ, was later reinterpreted in 1924 by de Broglie as a
standing wave condition. The wave-like properties of matter were subsequently confirmed by observations of electron interference
when scattered from crystals. Electrons can exist only in locations where they interfere constructively. How does this affect
electrons in atomic orbits? When an electron is bound to an atom, its wavelength must fit into a small space, something like a
standing wave on a string.
Waves on a String: (a) Waves on a string have a wavelength related to the length of the string, allowing them to interfere
constructively. (b) If we imagine the string bent into a closed circle, we get a rough idea of how electrons in circular orbits can
interfere constructively. (c) If the wavelength does not fit into the circumference, the electron interferes destructively; it cannot
exist in such an orbit.
Allowed orbits are those in which an electron constructively interferes with itself. Not all orbits produce constructive interference
and thus only certain orbits are allowed (i.e., the orbits are quantized). By assuming that the electron is described by a wave and a
whole number of wavelengths must fit along the circumference of the electron’s orbit, we have the equation:
nλ = 2πr (29.2.13)
Substituting de Broglie’s wavelength of hphp reproduces Bohr’s rule. Since λ=h/mevλ=h/mev, we now have:
nh
= 2π rn (29.2.14)
mev
Rearranging terms, and noting that L=mvrL=mvr for a circular orbit, we obtain the quantization of angular momentum as the
condition for allowed orbits:
h
L = me vrn = n , (n = 1, 2, 3 …) (29.2.15)
2π
As previously stated, Bohr was forced to hypothesize this rule for allowed orbits. We now realize this as the condition for
constructive interference of an electron in a circular orbit.
Accordingly, a new kind of mechanics, quantum mechanics, was proposed in 1925. Bohr’s model of electrons traveling in
quantized orbits was extended into a more accurate model of electron motion. The new theory was proposed by Werner
Heisenberg. By different reasoning, another form of the same theory, wave mechanics, was discovered independently by Austrian
physicist Erwin Schrödinger. Schrödinger employed de Broglie’s matter waves, but instead sought wave solutions of a three-
dimensional wave equation. This described electrons that were constrained to move about the nucleus of a hydrogen-like atom by
being trapped by the potential of the positive nuclear charge.
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de Broglie's Matter Waves Justify Bohr's Magic Elect…
Elect…
de Broglie’s Matter Waves Justify Bohr’s Magic Electron Orbital Radii: I include a summary of the hydrogen atom’s
electronic structure and explain how an electron can interfere with itself in an orbit just like it can in a double-slit experiment.
learning objectives
Describe Compton effects between electrons and x-ray photons
By the early 20th century, research into the interaction of X-rays with matter was well underway. It was observed that when X-rays
of a known wavelength interact with atoms, the X-rays are scattered through an angle θθ and emerge at a different wavelength
related to θθ. Although classical electromagnetism predicted that the wavelength of scattered rays should be equal to the initial
wavelength, multiple experiments had found that the wavelength of the scattered rays was longer (corresponding to lower energy)
than the initial wavelength.
In 1923, Compton published a paper in the Physical Review which explained the X-ray shift by attributing particle-like momentum
to “photons,” which Einstein had invoked in his Nobel prize winning explanation of the photoelectric effect. First postulated by
Planck, these “particles” conceptualized “quantized” elements of light as containing a specific amount of energy depending only on
the frequency of the light. In his paper, Compton derived the mathematical relationship between the shift in wavelength and the
scattering angle of the X-rays by assuming that each scattered X-ray photon interacted with only one electron. His paper concludes
by reporting on experiments which verified his derived relation:
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′
h
λ −λ = (1 − cos θ) (29.2.16)
me c
where λ\lambda is the initial wavelength, λ′\lambda’ is the wavelength after scattering, h is the Planck constant, m is the Electron
e
rest mass, c is the speed of light, and θ\theta is the scattering angle. The quantityh
me c
is known as the Compton wavelength of the
electron; it is equal to 2.43 ⋅ 10
−12
m. The wavelength shift λ‘−λ is at least zero (for θ = 0° ) and at most twice the Compton
wavelength of the electron (for θ = 0° ). (The derivation of Compton’s formula is a bit lengthy and will not be covered here.)
A Photon Colliding with a Target at Rest: A photon of wavelength λλ comes in from the left, collides with a target at rest, and a
new photon of wavelength λ‘λ‘ emerges at an angle θθ.
Because the mass-energy and momentum of a system must both be conserved, it is not generally possible for the electron simply to
move in the direction of the incident photon. The interaction between electrons and high energy photons (comparable to the rest
energy of the electron, 511 keV) results in the electron being given part of the energy (making it recoil), and a photon containing
the remaining energy being emitted in a different direction from the original, so that the overall momentum of the system is
conserved. If the scattered photon still has enough energy left, the Compton scattering process may be repeated. In this scenario,
the electron is treated as free or loosely bound. Photons with an energy of this order of magnitude are in the x-ray range of the
electromagnetic radiation spectrum. Therefore, you can say that Compton effects (with electrons) occur with x-ray photons.
If the photon is of lower energy, but still has sufficient energy (in general a few eV to a few keV, corresponding to visible light
through soft X-rays), it can eject an electron from its host atom entirely (a process known as the photoelectric effect), instead of
undergoing Compton scattering. Higher energy photons (1.022 MeV and above, in the gamma ray range) may be able to bombard
the nucleus and cause an electron and a positron to be formed, a process called pair production.
learning objectives
Describe interactions between X-rays and atoms
In a previous Atom on X-rays, we have seen that there are two processes by which x-rays are produced in the anode of an x-ray
tube. In one process, the deceleration of electrons produces x-rays, and these x-rays are called Bremsstrahlung, or braking
radiation. The second process is atomic in nature and produces characteristic x-rays, so called because they are characteristic of the
anode material. The x-ray spectrum in is typical of what is produced by an x-ray tube, showing a broad curve of Bremsstrahlung
radiation with characteristic x-ray peaks on it.
29.2.14 https://phys.libretexts.org/@go/page/16232
X-Ray Spectrum: X-ray spectrum obtained when energetic electrons strike a material, such as in the anode of a CRT. The smooth
part of the spectrum is bremsstrahlung radiation, while the peaks are characteristic of the anode material. A different anode material
would have characteristic x-ray peaks at different frequencies.
Since x-ray photons are very energetic, they have relatively short wavelengths. For example, the 54.4-keV Kα x-ray, for example,
has a wavelength λ = hc
E
= 0.0228 nm. Thus, typical x-ray photons act like rays when they encounter macroscopic objects, like
teeth, and produce sharp shadows. However, since atoms and atomic structures have a typical size on the order of 0.1 nm, x-ray
shows its wave nature with them. The process is called x-ray diffraction because it involves the diffraction and interference of x-
rays to produce patterns that can be analyzed for information about the structures that scattered the x-rays.
Shown below, Bragg’s Law gives the angles for coherent and incoherent scattering of light from a crystal lattice, which happens
during x-ray diffraction. When x-ray are incident on an atom, they make the electronic cloud move as an electromagnetic wave.
The movement of these charges re-radiate waves with the same frequency. This is called Rayleigh Scattering, which you should
remember from a previous atom. A similar thing happens when neutron waves from the nuclei scatter from interaction with an
unpaired electron. These re-emitted wave fields interfere with each other either constructively or destructively, and produce a
diffraction pattern that is captured by a sensor or film. This is called the Braggs diffraction, and is the basis for x-ray diffraction.
X-Ray Diffraction: Bragg’s Law of diffraction: illustration of how x-rays interact with crystal lattice.
Perhaps the most famous example of x-ray diffraction is the discovery of the double-helix structure of DNA in 1953. Using x-ray
diffraction data, researchers were able to discern the structure of DNA shows a diffraction pattern produced by the scattering of x-
rays from a crystal of protein. This process is known as x-ray crystallography because of the information it can yield about crystal
structure. Not only do x-rays confirm the size and shape of atoms, they also give information on the atomic arrangements in
materials. For example, current research in high-temperature superconductors involves complex materials whose lattice
arrangements are crucial to obtaining a superconducting material. These can be studied using x-ray crystallography.
X-Ray Diffraction: X-ray diffraction from the crystal of a protein, hen egg lysozyme, produced this interference pattern. Analysis
of the pattern yields information about the structure of the protein.
29.2.15 https://phys.libretexts.org/@go/page/16232
The Compton Effect
The Compton Effect is the phenomenon of the decrease in energy of photon when scattered by a free charged particle.
learning objectives
Explain why Compton scattering is an inelastic scattering.
Compton scattering is an inelastic scattering of a photon by a free charged particle (usually an electron). It results in a decrease in
energy (increase in wavelength) of the photon (which may be an X-ray or gamma ray photon), called the Compton Effect. Part of
the energy of the photon is transferred to the scattering electron. Inverse Compton scattering also exists, and happens when a
charged particle transfers part of its energy to a photon.
Scattering in the Compton Effect: The Compton Effect is the name given to the scattering of a photon by an electron. Energy and
momentum are conserved, resulting in a reduction of both for the scattered photon. Studying this effect, Compton verified that
photons have momentum.
Compton scattering is an example of inelastic scattering because the wavelength of the scattered light is different from the incident
radiation. Still, the origin of the effect can be considered as an elastic collision between a photon and an electron. The amount of
change in the wavelength is called the Compton shift. Although nuclear Compton scattering exists, Compton scattering usually
refers to the interaction involving only the electrons of an atom.
The Compton effect is important because it demonstrates that light cannot be explained purely as a wave phenomenon. Thomson
scattering, the classical theory of an electromagnetic wave scattered by charged particles, cannot explain low intensity shifts in
wavelength: classically, light of sufficient intensity for the electric field to accelerate a charged particle to a relativistic speed will
cause radiation-pressure recoil and an associated Doppler shift of the scattered light. However, the effect will become arbitrarily
small at sufficiently low light intensities regardless of wavelength. Light must behave as if it consists of particles to explain the
low-intensity Compton scattering. Compton’s experiment convinced physicists that light can behave as a stream of particle-like
objects (quanta) whose energy is proportional to the frequency.
Key Points
The British physicist J. J. Thomson performed experiments studying cathode rays and discovered that they were unique
particles, later named electrons.
Rutherford proved that the hydrogen nucleus is present in other nuclei.
In 1932, James Chadwick showed that there were uncharged particles in the radiation he was using. These particles, later called
neutrons, had a similar mass of the protons but did not have the same characteristics as protons.
The atom is a basic unit of matter that consists of a dense central nucleus surrounded by a cloud of negatively charged electrons.
Scattered knowledge discovered by alchemists over the Middle Ages contributed to the discovery of atoms.
Dalton established his atomic theory based on the fact that the masses of reactants in specific chemical reactions always have a
particular mass ratio.
J. J. Thomson, who discovered the electron in 1897, proposed the plum pudding model of the atom in 1904 before the discovery
of the atomic nucleus in order to include the electron in the atomic model.
In Thomson’s model, the atom is composed of electrons surrounded by a soup of positive charge to balance the electrons’
negative charges, like negatively charged “plums” surrounded by positively charged “pudding”.
The 1904 Thomson model was disproved by Hans Geiger’s and Ernest Marsden’s 1909 gold foil experiment.
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Rutherford overturned Thomson’s model in 1911 with his well-known gold foil experiment, in which he demonstrated that the
atom has a tiny, high- mass nucleus.
In his experiment, Rutherford observed that many alpha particles were deflected at small angles while others were reflected
back to the alpha source.
This highly concentrated, positively charged region is named the “nucleus” of the atom.
According to Bohr: 1) Electrons in atoms orbit the nucleus, 2) The electrons can only orbit stably, without radiating, in certain
orbits, and 3) Electrons can only gain and lose energy by jumping from one allowed orbit to another.
The significance of the Bohr model is that the laws of classical mechanics apply to the motion of the electron about the nucleus
only when restricted by a quantum rule. Therefore, his atomic model is called a semiclassical model.
The laws of classical mechanics predict that the electron should release electromagnetic radiation while orbiting a nucleus,
suggesting that all atoms should be unstable!
Classical electrodynamics predicts that an atom described by a (classical) planetary model would be unstable.
To explain the hydrogen spectrum, Bohr had to make a few assumptions that electrons could only have certain classical
motions.
After the seminal work by Planck, Einstein, and Bohr, physicists began to realize that it was essential to introduce the notion of
” quantization ” to explain microscopic worlds.
The “Bohr orbits” have a very important feature of quantization: that the angular momentum L of an electron in its orbit is
quantized, that is, it has only specific, discrete values. This leads to the equation L = m vr = n = nℏ .
e n
h
2π
At the time of proposal, Bohr himself did not know why angular momentum should be quantized, but using this assumption he
was able to calculate the energies in the hydrogen spectrum.
A theory of the atom or any other system must predict its energies based on the physics of the system, which the Bohr model
was able to do.
2 2
The energy of a photon emitted by a hydrogen atom is given by the difference of two hydrogen energy levels:
E = Ei − Ef = R (
n
1
2
−
1
2
n
) , which is known as Rydberg formula.
f i
Atomic and molecular emission and absorption spectra have been known for over a century to be discrete (or quantized).
Lyman, Balmer, and Paschen series are named after early researchers who studied them in particular depth.
Bohr was the first one to provide a theoretical explanation of the hydrogen spectra.
Bohr’s condition, that the angular momentum is an integer multiple of ℏℏ, was later reinterpreted in 1924 by de Broglie as a
standing wave condition.
For what Bohr was forced to hypothesize as the rule for allowed orbits, de Broglie’s matter wave concept explains it as the
condition for constructive interference of an electron in a circular orbit.
Bohr’s model was only applicable to hydrogen-like atoms. In 1925, more general forms of description (now called quantum
mechanics ) emerged, thanks to Heisenberg and Schrodinger.
Compton derived the mathematical relationship between the shift in wavelength and the scattering angle of the X-rays.
Compton effects (with electrons) usually occur with x-ray photons.
If the photon is of lower energy, in the visible light through soft X-rays range, photoelectric effects are observed. Higher energy
photons, in the gamma ray range, may lead to pair production.
X rays are relatively high- frequency EM radiation. They are produced by transitions between inner-shell electron levels, which
produce x rays characteristic of the atomic element, or by accelerating electrons.
x-ray diffraction is a technique that provides the detailed information about crystallographic structure of natural and
manufactured materials.
Current research in material science and physics involves complex materials whose lattice arrangements are crucial to obtaining
a superconducting material, which can be studied using x-ray crystallography.
Key Terms
scintillation: A flash of light produced in a transparent material by the passage of a particle.
alpha particle: A positively charged nucleus of a helium-4 atom (consisting of two protons and two neutrons), emitted as a
consequence of radioactivity.
29.2.17 https://phys.libretexts.org/@go/page/16232
cathode: An electrode through which electric current flows out of a polarized electrical device.
electromagnetic force: a long-range fundamental force that acts between charged bodies, mediated by the exchange of photons
Avogadro’s number: the number of constituent particles (usually atoms or molecules) in one mole of a given substance. It has
dimensions of reciprocal mol and its value is equal to 6.02214129 ⋅ 10 mol
23 −1
nucleus: the massive, positively charged central part of an atom, made up of protons and neutrons
Maxwell’s equations: A set of equations describing how electric and magnetic fields are generated and altered by each other
and by charges and currents.
semiclassical: a theory in which one part of a system is described quantum-mechanically whereas the other is treated
classically.
black body: An idealized physical body that absorbs all incident electromagnetic radiation, regardless of frequency or angle of
incidence. Although black body is a theoretical concept, you can find approximate realizations of black body in nature.
photoelectric effect: The occurrence of electrons being emitted from matter (metals and non-metallic solids, liquids, or gases)
as a consequence of their absorption of energy from electromagnetic radiation.
quantization: The process of explaining a classical understanding of physical phenomena in terms of a newer understanding
known as quantum mechanics.
centripetal: Directed or moving towards a center.
photon: The quantum of light and other electromagnetic energy, regarded as a discrete particle having zero rest mass, no
electric charge, and an indefinitely long lifetime.
spectrum: A condition that is not limited to a specific set of values but can vary infinitely within a continuum. The word saw
its first scientific use within the field of optics to describe the rainbow of colors in visible light when separated using a prism.
standing wave: A wave form which occurs in a limited, fixed medium in such a way that the reflected wave coincides with the
produced wave. A common example is the vibration of the strings on a musical stringed instrument.
matter wave: A concept reflects the wave-particle duality of matter. The theory was proposed by Louis de Broglie.
gamma ray: A very high frequency (and therefore very high energy) electromagnetic radiation emitted as a consequence of
radioactivity.
photoelectric effects: In photoelectric effects, electrons are emitted from matter (metals and non-metallic solids, liquids or
gases) as a consequence of their absorption of energy from electromagnetic radiation.
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29.3: Atomic Physics and Quantum Mechanics
learning objectives
Explain relationship between the wave nature of matter and the quantization of energy levels in bound systems
To consider why wave nature of matter in bound systems leads to quantization, let’s consider an example in classical mechanics.
We will look at a basic string “instrument” (a string pulled tight and fixed at both ends). If a string was free and not attached to
anything, we know that it could oscillate at any driven frequency. However, the string in this example (with fixed ends and specific
length) can only produce a very specific set of pitches because only waves of a certain wavelength can “fit” on the string of a given
length with fixed ends. Once the string becomes a “bound system” with specific boundary restrictions, it allows waves with only a
discrete set of frequencies.
This is the exact mechanism that causes quantization in atoms. The wave nature of matter is responsible for the quantization of
energy levels in bound systems. Just like a free string, the matter wave of a free electron can have any wavelength, determined by
its momentum. However, once an electron is “bound” by a Coulomb potential of a nucleus, it can no longer have an arbitrary
wavelength as the wave needs to satisfy a certain boundary condition. Only those states where matter interferes constructively
(leading to standing waves) exist, or are “allowed” (see illustration in.
Fig 2: The third and fourth allowed circular orbits have three and four wavelengths, respectively, in their circumferences.
Assuming that an integral multiple of the electron’s wavelength equals the circumference of the orbit, we have:
nλn = 2π rn (n = 1, 2, 3, …) (29.3.1)
Substituting λ = h
me v
, this becomes:
nh
= 2π rn (29.3.2)
mev
As previously discussed, Bohr was forced to hypothesize this as the rule for allowed orbits. We now realize this as a condition for
constructive interference of an electron in a (bound) circular orbit.
learning objectives
Describe process of pair production as the result of photon interaction with nucleus
Below is an illustration of pair production, which refers to the creation of an elementary particle and its antiparticle, usually when a
photon interacts with a nucleus. For example, an electron and its antiparticle, the positron, may be created. This is allowed,
provided there is enough energy available to create the pair (i.e., the total rest mass energy of the two particles) and that the
29.3.1 https://phys.libretexts.org/@go/page/16234
situation allows both energy and momentum to be conserved. Some other conserved quantum numbers such as angular momentum,
electric charge, etc., must sum to zero as well. The probability of pair production in photon-matter interactions increases with
increasing photon energy, and also increases with atomic number (Z) of the nucleus approximately as Z .
2
Pair Production: Feynman diagram for pair production. A photon decays into an electron-positron pair.
− +
γ + γ → e + e
In nuclear physics, this reaction occurs when a high-energy photon ( gamma rays ) interacts with a nucleus. The energy of this
photon can be converted into mass through Einstein’s equation E = mc where E is energy, mm is mass and c is the speed of light.
2
The photon must have enough energy to create the mass of an electron plus a positron. The mass of an electron is 9.11 ⋅ 10 kg
−31
Key Points
Strings in musical instruments (guitar, for example) can only produce a very specific set of pitches because only waves of a
certain wavelength can “fit” on the string of a given length with fixed ends.
Similarly, once an electron is bound by a Coulomb potential of a nucleus, it no longer can have any arbitrary wavelength
because the wave should satisfy a certain boundary condition.
Bohr’s quantization assumption can be derived from the condition for constructive intereference of an electron matter wave in a
circular orbit.
The probability of pair production in photon -matter interactions increases with increasing photon energy, and also increases
with atomic number of the nucleus approximately as Z.
Energy and momentum should be conserved through the pair production process. Some other conserved quantum numbers such
as angular momentum, electric charge, etc., must sum to zero as well.
Nucleus is needed in the pair production of electron and positron to satisfy the energy and momentum conservation laws.
Key Terms
quantization: The process of explaining a classical understanding of physical phenomena in terms of a newer understanding
known as quantum mechanics.
angular momentum: A vector quantity describing an object in circular motion; its magnitude is equal to the momentum of the
particle, and the direction is perpendicular to the plane of its circular motion.
matter wave: A concept reflects the wave-particle duality of matter. The theory was proposed by Louis de Broglie.
gamma ray: A very high frequency (and therefore very high energy) electromagnetic radiation emitted as a consequence of
radioactivity.
positron: The antimatter equivalent of an electron, having the same mass but a positive charge.
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29.4: Applications of Atomic Physics
We have seen that under certain circumstances particles behave like waves. This idea is used in the electron microscope which is a
type that uses electrons to create an image of the target. It has much higher magnification or resolving power than a normal light
microscope. It can achieve better than 50 pm resolution and magnifications of up to about 10,000,000 times, whereas ordinary,
nonconfocal light microscopes are limited by diffraction to about 200 nm resolution and useful magnifications below 2000 times.
Optical and Electron Microscopes: Diagram of the basic components of an optical microscope and an electron microscope.
The original form of electron microscopy, transmission electron microscopy, works in a similar manner using electrons. In the
electron microscope, electrons which are emitted by a cathode are formed into a beam using magnetic lenses (usually
electromagnets). This electron beam is then passed through a very thin target. Again, the regions in the target with higher densities
stop the electrons more easily. So, the number of electrons which pass through the different regions of the target depends on their
densities. This means that the partially transmitted beam of electrons carries information about the densities of the inner structure of
the target.
The spatial variation in this information (the “image”) is then magnified by a series of magnetic lenses and it is recorded by hitting
a fluorescent screen, photographic plate, or light-sensitive sensor such as a CCD (charge-coupled device) camera. The image
detected by the CCD may be displayed in real time on a monitor or computer.
Electron microscopes are very useful as they are able to magnify objects to a much higher resolution. This is because their de
Broglie wavelengths are so much smaller than that of visible light. You hopefully remember that light is diffracted by objects which
are separated by a distance of about the same size as the wavelength of the light. This diffraction then prevents you from being able
to focus the transmitted light into an image.
Therefore, the sizes at which diffraction occurs for a beam of electrons is much smaller than those for visible light. This is why you
can magnify targets to a much higher order of magnification using electrons rather than visible light.
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Lasers
A laser consists of a gain medium, a mechanism to supply energy to it, and something to provide optical feedback.
learning objectives
Describe basic parts of laser
When lasers were invented in 1960, they were called “a solution looking for a problem. ” Nowadays, lasers are ubiquitous, finding
utility in thousands of highly varied applications in every section of modern society, including consumer electronics, information
technology, science, medicine, industry, law enforcement, entertainment, and the military.
Having examined stimulated emission and optical amplification process in the “Lasers, Applications of Quantum Mechanics”
section, this atom looks at how lasers are built.
A laser consists of a gain medium, a mechanism to supply energy to it, and something to provide optical feedback (usually an
optical cavity). When a gain medium is placed in an optical cavity, a laser can then produce a coherent beam of photons.
The gain medium is where the optical amplification process occurs. It is excited by an external source of energy into an excited
state (called “population inversion”), ready to be fired when a photon with the right frequency enters the medium. In most lasers,
this medium consists of a population of atoms which have been excited by an outside light source or an electrical field which
supplies energy for atoms to absorb in order to be transformed into excited states. There are many types of lasers depending on the
gain media and mode of operation. Gas and semiconductors are commonly used gain media.
Ruby (doubled) 347 nm
Er:glass 1.54 µm
Nd:YAG 1319 nm
KrCl excimer 222 nm
Xe 3+ 539.5 nm
Nd:YAG 1064 nm
CO 2 10.6 µm
Ruby 694.3 nm
Nd:YAG (tripled) 355 nm
Pulse CW
Nitrogen 337.1 nm
F2 excimer 157 nm
Er:YSGG 2.79 µm
Ar + 488.0 nm
Kr + 647.1 nm
Tm:YAG 2.01 µm
Ar + 514.5 nm
Ho:YAG 2.08 µm
energy power
He-Ag + 224.3 nm
Nd:YAG 946 nm
He-Ne 3.391 µm
He-Cd 441.6 nm
He-Ne 543.5 nm
He-Ne 594.1 nm
nm
He-Ne 611.9 nm
He-Ne 1152 nm
He-Cd 325 nm
He-Ne 632.8
1 kJ 1 kW
1J 1W
1 mJ 1 mW X-RAYS FAR-IR
ULTRAVIOLET VISIBLE NEAR-INFRARED MID-INFRARED
100 nm 200 nm 300 nm 400 nm 500 nm 600 nm 700 nm 800 nm 900 nm 1 µm 3 µm 10 µm 30 µm 11 mm
1 mJ 1 mW Dyes in
polymer InGaAs Xe-He
550-700 nm 904-1065 nm 2-4 µm
Dyes CO 2 (doubled) Pb salts
Ti:sapphire 4.6-5.8 µm
1J 1W He-Au
+
0.38-1.0 µm 3.3-27 µm
(tripled) Ti:sapphire InGaAlP
Ne-Cu 282-292 nm
+
235-330 nm (doubled) Cr fluoride CO
Alexandrite 630-685 nm CO 2
248-270 nm 360-460 nm 780-850 nm 5-7 µm
(doubled) Alexandrite InGaAs 9.2-11.4 µm
1 kJ 1 kW Dyes (doubled) GaN DF chemical
360-400 nm 700-800 nm 1.27-1.33 µm
0.2-0.4 µm 515-520 nm 3.6-4.0 µm
InGaN GaAlAs 1.43-1.57 µm
Ti:sapphire Fosterite HF chemical
370-493 nm 750-850 nm
670-1130 nm 1.13-1.36 µm 2.6-3.0 µm
AlGaIn/AsSb
1.87-2.2 µm
Dyes (Raman shifted)
0.9-4.5 µm
Wavelengths of Commercially Available Lasers: Laser types with distinct laser lines are shown above the wavelength bar, while
below are shown lasers that can emit in a wavelength range. The height of the lines and bars gives an indication of the maximal
power/pulse energy commercially available, while the color codifies the type of laser material.
The most common type of laser uses feedback from an optical cavity–a pair of highly reflective mirrors on either end of the gain
medium. A single photon can bounce back and forth between the mirrors many times, passing through the gain medium and being
amplified each time. Typically one of the two mirrors, the output coupler, is partially transparent. Some of the light escapes through
this mirror, producing a laser beam that is visible to the naked eye.
Key Points
Electron microscopes are very useful as they are able to magnify objects to a much higher resolution than optical ones.
Higher resolution can be achieved with electron microscopes because the de Broglie wavelengths for electrons are so much
smaller than that of visible light.
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In electron microscopes, electromagnets can be used as magnetic lenses to manipulate electron beams.
The gain medium is where the optical amplification process occurs. Gas and semiconductors are commonly used gain media.
The most common type of laser uses feedback from an optical cavity–a pair of highly reflective mirrors on either end of the
gain medium. A single photon can bounce back and forth between the mirrors many times, passing through the gain medium
and being amplified each time.
Lasers are ubiquitous, finding utility in thousands of highly varied applications in every section of modern society.
Key Terms
CCD: A charge-coupled device (CCD) is a device for the movement of electrical charge, usually from within the device to an
area where the charge can be manipulated, for example conversion into a digital value. The CCD is a major technology required
for digital imaging.
de Broglie wavelength: The wavelength of a matter wave is inversely proportional to the momentum of a particle and is called
a de Broglie wavelength.
stimulated emission: The process by which an atomic electron (or an excited molecular state) interacting with an
electromagnetic wave of a certain frequency may drop to a lower energy level, transferring its energy to that field.
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29.5: Multielectron Atoms
learning objectives
Describe atomic structure and shielding in multielectron atoms
Multielectron Atoms
Atoms with more than one electron, such as Helium (He) and Nitrogen (N), are referred to as multielectron atoms. Hydrogen is the
only atom in the periodic table that has one electron in the orbitals under ground state.
In hydrogen-like atoms (those with only one electron), the net force on the electron is just as large as the electric attraction from the
nucleus. However, when more electrons are involved, each electron (in the nn-shell) feels not only the electromagnetic attraction
from the positive nucleus, but also repulsion forces from other electrons in shells from ‘1’ to ‘nn‘. This causes the net force on
electrons in the outer electron shells to be significantly smaller in magnitude. Therefore, these electrons are not as strongly bonded
to the nucleus as electrons closer to the nucleus. This phenomenon is often referred to as the Orbital Penetration Effect. The
shielding theory also explains why valence shell electrons are more easily removed from the atom.
Electron Shielding Effect: A multielectron atom with inner electrons shielding outside electrons from the positively charged
nucleus
The size of the shielding effect is difficult to calculate precisely due to effects from quantum mechanics. As an approximation, the
effective nuclear charge on each electron can be estimated by: Z = Z − σZ = Z– σ , where Z is the number of protons in the
ef f eff
nucleus and σ\sigma is the average number of electrons between the nucleus and the electron in question. σ\sigma can be found by
using quantum chemistry and the Schrodinger equation or by using Slater’s empirical formula.
For example, consider a sodium cation, a fluorine anion, and a neutral neon atom. Each has 10 electrons, and the number of
nonvalence electrons is two (10 total electrons minus eight valence electrons), but the effective nuclear charge varies because each
has a different number of protons:
−
Zef f F = 9 − 2 = 7+ (29.5.1)
Zef f Ne = 10 − 2 = 8+ (29.5.2)
+
Zef f Na = 11 − 2 = 9+ (29.5.3)
As a consequence, the sodium cation has the largest effective nuclear charge and, therefore, the smallest atomic radius.
learning objectives
Explain how elements are arranged in the Periodic Table.
A periodic table is a tabular display of the chemical elements, organized on the basis of their atomic numbers, electron
configurations, and recurring chemical properties. Elements are presented according to their atomic numbers (number of protons)
in increasing order. The standard form of the table comprises an eighteen by seven grid or main body of elements, positioned above
a smaller double row of elements. The table can also be deconstructed into four rectangular blocks: the s-block to the left, the p-
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block to the right, the d-block in the middle, and the f-block below that. The rows of the table are called periods. The columns of
the s-, d-, and p-blocks are called groups, some of which have names such as the halogens or the noble gases.
Since, by definition, a periodic table incorporates recurring trends, any such table can be used to derive relationships between the
properties of the elements and predict the properties of new elements that are yet to be discovered or synthesized. As a result, a
periodic table, in the standard form or some other variant, provides a useful framework for analyzing chemical behavior. Such
tables are widely used in chemistry and other sciences.
Periodic Table of Elements: The standard form of the periodic table, where the colors represent different categories of elements
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Mendeleev’s 1869 Periodic Table: Mendeleev’s 1869 periodic table presents the periods vertically and the groups horizontally.
Dmitri Mendeleev: Dmitri Mendeleev is known for publishing a widely recognized periodic table.
Electron Configurations
The electron configuration is the distribution of electrons of an atom or molecule in atomic or molecular orbitals.
learning objectives
Explain the meaning of electron configurations
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The electron configuration is the distribution of electrons of an atom or molecule in atomic or molecular orbitals. Electron
configurations describe electrons as each moving independently in an orbital, in an average field created by all other orbitals.
In atoms, electrons fill atomic orbitals according to the Aufbau principle (shown in ), stated as: a maximum of two electrons are put
into orbitals in the order of increasing orbital energy—the lowest-energy orbitals are filled before electrons are placed in higher-
energy orbitals. As an example, the electron configuration of the neon atom is 1s2 2s2 2p6 or [He]2s2 2p6, as diagramed in. In
molecules, the situation becomes more complex, as each molecule has a different orbital structure. The molecular orbitals are
labelled according to their symmetry, rather than the atomic orbital labels used for atoms and monoatomic ions: hence, the electron
configuration of the diatomic oxygen molecule, O2, is 1σg2 1σu2 2σg2 2σu2 1πu4 3σg2 1πg2.
Electron Configuration of Neon Atom: Electron configuration of neon atom showing only outer electron shell.
Aufbau Principle: In the Aufbau Principle, as electrons are added to atoms, they are added to the lowest orbitals first.
According to the laws of quantum mechanics, for systems with only one electron, an energy is associated with each electron
configuration and, upon certain conditions, electrons are able to move from one configuration to another by emission or absorption
of a quantum of energy, in the form of a photon.
For atoms or molecules with more than one electron, the motion of electrons are correlated and such picture is no longer exact. An
infinite number of electronic configurations are needed to exactly describe any multi-electron system, and no energy can be
associated with one single configuration. However, the electronic wave function is usually dominated by a very small number of
configurations and therefore the notion of electronic configuration remains essential for multi-electron systems.
Electronic configuration of polyatomic molecules can change without absorption or emission of photon through vibronic couplings.
Knowledge of the electron configuration of different atoms is useful in understanding the structure of the periodic table of
elements. The outermost electron shell is often referred to as the valence shell and (to a first approximation) determines the
chemical properties. It should be remembered that the similarities in the chemical properties were remarked more than a century
before the idea of electron configuration. The concept of electron configuration is also useful for describing the chemical bonds that
hold atoms together. In bulk materials this same idea helps explain the peculiar properties of lasers and semiconductors.
Key Points
Hydrogen is the only atom in the periodic table that has one electron in the orbitals under ground state.
In multielectron atoms, the net force on electrons in the outer shells is reduced due to shielding.
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The effective nuclear charge on each electron can be approximated as: Zeff=Z−σZeff=Z−σ, where ZZ is the number of protons
in the nucleus and σσ is the average number of electrons between the nucleus and the electron in question.
A periodic table provides a useful framework for analyzing the chemical behavior of elements.
A periodic table includes only chemical elements with each chemical element assigned a unique atomic number representing
the number of protons in its nucleus.
Dmitri Mendeleev is credited with the publication of the first widely recognized periodic table in 1869.
Electrons fill atomic orbitals according to the Aufbau principle in atoms.
For systems with only one electron, an energy is associated with each electron configuration and electrons are able to move
from one configuration to another by emission or absorption of a quantum of energy, in the form of a photon.
For atoms or molecules with more than one electron, an infinite number of electronic configurations are needed to exactly
describe any multi-electron system, and no energy can be associated with one single configuration.
Key Terms
hydrogen-like: having a single electron
electron shell: The collective states of all electrons in an atom having the same principal quantum number (visualized as an
orbit in which the electrons move).
valence shell: the outermost shell of electrons in an atom; these electrons take part in bonding with other atoms
periodic table: A tabular chart of the chemical elements according to their atomic numbers so that elements with similar
properties are in the same column.
element: Any one of the simplest chemical substances that cannot be decomposed in a chemical reaction or by any chemical
means and made up of atoms all having the same number of protons.
atomic number: The number, equal to the number of protons in an atom that determines its chemical properties. Symbol: Z
atomic orbital: The quantum mechanical behavior of an electron in an atom describing the probability of the electron’s
particular position and energy.
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CHAPTER OVERVIEW
3: Two-Dimensional Kinematics
3.1: Motion in Two Dimensions
3.2: Vectors
3.3: Projectile Motion
3.4: Multiple Velocities
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1
3.1: Motion in Two Dimensions
Constant Velocity
An object moving with constant velocity must have a constant speed in a constant direction.
learning objectives
Examine the terms for constant velocity and how they apply to acceleration
Motion with constant velocity is one of the simplest forms of motion. This type of motion occurs when an an object is moving (or
sliding) in the presence of little or negligible friction, similar to that of a hockey puck sliding across the ice. To have a constant
velocity, an object must have a constant speed in a constant direction. Constant direction constrains the object to motion to a
straight path.
Newton’s second law (F = ma ) suggests that when a force is applied to an object, the object would experience acceleration. If the
acceleration is 0, the object shouldn’t have any external forces applied on it. Mathematically, this can be shown as the following:
dv
a = = 0 ⇒ v = const. (3.1.1)
dt
If an object is moving at constant velocity, the graph of distance vs. time (x vs. t ) shows the same change in position over each
interval of time. Therefore the motion of an object at constant velocity is represented by a straight line: x = x + vt , where x is
0 0
Motion with Constant Velocity: When an object is moving with constant velocity, it does not change direction nor speed and
therefore is represented as a straight line when graphed as distance over time.
You can also obtain an object’s velocity if you know its trace over time. Given a graph as in, we can calculate the velocity from the
change in distance over the change in time. In graphical terms, the velocity can be interpreted as the slope of the line. The velocity
can be positive or negative, and is indicated by the sign of our slope. This tells us in which direction the object moves.
Constant Acceleration
Analyzing two-dimensional projectile motion is done by breaking it into two motions: along the horizontal and vertical axes.
learning objectives
Analyze a two-dimensional projectile motion along horizontal and vertical axes
Projectile motion is the motion of an object thrown, or projected, into the air, subject only to the force of gravity. The object is
called a projectile, and its path is called its trajectory. The motion of falling objects is a simple one-dimensional type of projectile
motion in which there is no horizontal movement. In two-dimensional projectile motion, such as that of a football or other thrown
object, there is both a vertical and a horizontal component to the motion.
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Projectile Motion: Throwing a rock or kicking a ball generally produces a projectile pattern of motion that has both a vertical and
a horizontal component.
The most important fact to remember is that motion along perpendicular axes are independent and thus can be analyzed separately.
The key to analyzing two-dimensional projectile motion is to break it into two motions, one along the horizontal axis and the other
along the vertical. To describe motion we must deal with velocity and acceleration, as well as with displacement.
We will assume all forces except for gravity (such as air resistance and friction, for example) are negligible. The components of
acceleration are then very simple: a = −g = −9.81
y (we assume that the motion occurs at small enough heights near the
m
2
s
surface of the earth so that the acceleration due to gravity is constant). Because the acceleration due to gravity is along the vertical
direction only, a = 0 . Thus, the kinematic equations describing the motion along the x and y directions respectively, can be used:
x
x = x0 + vx t (3.1.2)
y = v0y + ay t (3.1.3)
1 2
y = y0 + v0y t + ay t (3.1.4)
2
2 2
vy =v + 2 ay (y − y0 ) (3.1.5)
0y
We analyze two-dimensional projectile motion by breaking it into two independent one-dimensional motions along the vertical and
horizontal axes. The horizontal motion is simple, because a = 0 and v is thus constant. The velocity in the vertical direction
x x
begins to decrease as an object rises; at its highest point, the vertical velocity is zero. As an object falls towards the Earth again, the
vertical velocity increases again in magnitude but points in the opposite direction to the initial vertical velocity. The xx and yy
motions can be recombined to give the total velocity at any given point on the trajectory.
Key Points
Constant velocity means that the object in motion is moving in a straight line at a constant speed.
This line can be represented algebraically as: x = x + vt , where x represents the position of the object at t = 0 , and the
0 0
Key Terms
constant velocity: Motion that does not change in speed nor direction.
kinematic: of or relating to motion or kinematics
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3.2: Vectors
Components of a Vector
Vectors are geometric representations of magnitude and direction and can be expressed as arrows in two or three dimensions.
learning objectives
Contrast two-dimensional and three-dimensional vectors
Vectors are geometric representations of magnitude and direction which are often represented by straight arrows, starting at one
point on a coordinate axis and ending at a different point. All vectors have a length, called the magnitude, which represents some
quality of interest so that the vector may be compared to another vector. Vectors, being arrows, also have a direction. This
differentiates them from scalars, which are mere numbers without a direction.
A vector is defined by its magnitude and its orientation with respect to a set of coordinates. It is often useful in analyzing vectors to
break them into their component parts. For two-dimensional vectors, these components are horizontal and vertical. For three
dimensional vectors, the magnitude component is the same, but the direction component is expressed in terms of xx, yy and zz.
Decomposing a Vector
To visualize the process of decomposing a vector into its components, begin by drawing the vector from the origin of a set of
coordinates. Next, draw a straight line from the origin along the x-axis until the line is even with the tip of the original vector. This
is the horizontal component of the vector. To find the vertical component, draw a line straight up from the end of the horizontal
vector until you reach the tip of the original vector. You should find you have a right triangle such that the original vector is the
hypotenuse.
Decomposing a vector into horizontal and vertical components is a very useful technique in understanding physics problems.
Whenever you see motion at an angle, you should think of it as moving horizontally and vertically at the same time. Simplifying
vectors in this way can speed calculations and help to keep track of the motion of objects.
3.2.1 https://phys.libretexts.org/@go/page/14446
Scalars and Vectors: Mr. Andersen explains the differences between scalar and vectors quantities. He also uses a demonstration to
show the importance of vectors and vector addition.
Components of a Vector: The original vector, defined relative to a set of axes. The horizontal component stretches from the start
of the vector to its furthest x-coordinate. The vertical component stretches from the x-axis to the most vertical point on the vector.
Together, the two components and the vector form a right triangle.
learning objectives
Distinguish the difference between the quantities scalars and vectors represent
Physical quantities can usually be placed into two categories, vectors and scalars. These two categories are typified by what
information they require. Vectors require two pieces of information: the magnitude and direction. In contrast, scalars require only
the magnitude. Scalars can be thought of as numbers, whereas vectors must be thought of more like arrows pointing in a specific
direction.
3.2.2 https://phys.libretexts.org/@go/page/14446
A Vector: An example of a vector. Vectors are usually represented by arrows with their length representing the magnitude and their
direction represented by the direction the arrow points.
Vectors require both a magnitude and a direction. The magnitude of a vector is a number for comparing one vector to another. In
the geometric interpretation of a vector the vector is represented by an arrow. The arrow has two parts that define it. The two parts
are its length which represents the magnitude and its direction with respect to some set of coordinate axes. The greater the
magnitude, the longer the arrow. Physical concepts such as displacement, velocity, and acceleration are all examples of quantities
that can be represented by vectors. Each of these quantities has both a magnitude (how far or how fast) and a direction. In order to
specify a direction, there must be something to which the direction is relative. Typically this reference point is a set of coordinate
axes like the x-y plane.
Scalars differ from vectors in that they do not have a direction. Scalars are used primarily to represent physical quantities for which
a direction does not make sense. Some examples of these are: mass, height, length, volume, and area. Talking about the direction of
these quantities has no meaning and so they cannot be expressed as vectors.
3.2.3 https://phys.libretexts.org/@go/page/14446
The difference between Vectors and Scalars, Introduction and Basics: This video introduces the difference between scalars and
vectors. Ideas about magnitude and direction are introduced and examples of both vectors and scalars are given.
learning objectives
Distinguish the difference between the quantities scalars and vectors represent
Graphical Addition of Vectors: The head-to-tail method of vector addition requires that you lay out the first vector along a set of
coordinate axes. Next, place the tail of the next vector on the head of the first one. Draw a new vector from the origin to the head of
the last vector. This new vector is the sum of the original two.
3.2.4 https://phys.libretexts.org/@go/page/14446
Vector Addition Lesson 1 of 2: Head to Tail Addition …
Vector Addition Lesson 1 of 2: Head to Tail Addition Method: This video gets viewers started with vector addition and
subtraction. The first lesson shows graphical addition while the second video takes a more mathematical approach and shows
vector addition by components.
To subtract vectors the method is similar. Make sure that the first vector you draw is the one to be subtracted from. Then, to
subtract a vector, proceed as if adding the opposite of that vector. In other words, flip the vector to be subtracted across the axes and
then join it tail to head as if adding. To flip the vector, simply put its head where its tail was and its tail where its head was.
learning objectives
Demonstrate how to add and subtract vectors by components
3.2.5 https://phys.libretexts.org/@go/page/14446
components and then add the components together.
Vector with Horizontal and Vertical Components: The vector in this image has a magnitude of 10.3 units and a direction of 29.1
degrees above the x-axis. It can be decomposed into a horizontal part and a vertical part as shown.
For example, a vector with a length of 5 at a 36.9 degree angle to the horizontal axis will have a horizontal component of 4 units
and a vertical component of 3 units. If we were to add this to another vector of the same magnitude and direction, we would get a
vector twice as long at the same angle. This can be seen by adding the horizontal components of the two vectors (4 + 4) and the
two vertical components (3 + 3 ). These additions give a new vector with a horizontal component of 8(4 + 4) and a vertical
component of 6(3 + 3) . To find the resultant vector, simply place the tail of the vertical component at the head (arrow side) of the
horizontal component and then draw a line from the origin to the head of the vertical component. This new line is the resultant
vector. It should be twice as long as the original, since both of its components are twice as large as they were previously.
To subtract vectors by components, simply subtract the two horizontal components from each other and do the same for the vertical
components. Then draw the resultant vector as you did in the previous part.
3.2.6 https://phys.libretexts.org/@go/page/14446
Vector Addition Lesson 2 of 2: How to Add Vectors by Components: This video gets viewers started with vector addition using
a mathematical approach and shows vector addition by components.
learning objectives
Summarize the interaction between vectors and scalars
Overview
Although vectors and scalars represent different types of physical quantities, it is sometimes necessary for them to interact. While
adding a scalar to a vector is impossible because of their different dimensions in space, it is possible to multiply a vector by a
scalar. A scalar, however, cannot be multiplied by a vector.
To multiply a vector by a scalar, simply multiply the similar components, that is, the vector’s magnitude by the scalar’s magnitude.
This will result in a new vector with the same direction but the product of the two magnitudes.
Example 3.2.1:
For example, if you have a vector A with a certain magnitude and direction, multiplying it by a scalar a with magnitude 0.5
will give a new vector with a magnitude of half the original. Similarly if you take the number 3 which is a pure and unit-less
scalar and multiply it to a vector, you get a version of the original vector which is 3 times as long. As a more physical example
take the gravitational force on an object. The force is a vector with its magnitude depending on the scalar known as mass and
its direction being down. If the mass of the object is doubled, the force of gravity is doubled as well.
Multiplying vectors by scalars is very useful in physics. Most of the units used in vector quantities are intrinsically scalars
multiplied by the vector. For example, the unit of meters per second used in velocity, which is a vector, is made up of two
scalars, which are magnitudes: the scalar of length in meters and the scalar of time in seconds. In order to make this conversion
from magnitudes to velocity, one must multiply the unit vector in a particular direction by these scalars.
3.2.7 https://phys.libretexts.org/@go/page/14446
Scalar Multiplication: (i) Multiplying the vector A by the scalar a = 0.5 yields the vector B which is half as long. (ii)
Multiplying the vector A by 3 triples its length. (iii) Doubling the mass (scalar) doubles the force (vector) of gravity.
learning objectives
Predict the influence of multiplying a vector by a scalar
In addition to adding vectors, vectors can also be multiplied by constants known as scalars. Scalars are distinct from vectors in that
they are represented by a magnitude but no direction. Examples of scalars include an object’s mass, height, or volume.
Scalar Multiplication: (i) Multiplying the vector A by the scalar a = 0.5 yields the vector B which is half as long. (ii)
Multiplying the vector A by 3 triples its length. (iii) Doubling the mass (scalar) doubles the force (vector) of gravity.
When multiplying a vector by a scalar, the direction of the vector is unchanged and the magnitude is multiplied by the magnitude of
the scalar. This results in a new vector arrow pointing in the same direction as the old one but with a longer or shorter length. You
can also accomplish scalar multiplication through the use of a vector’s components. Once you have the vector’s components,
multiply each of the components by the scalar to get the new components and thus the new vector.
A useful concept in the study of vectors and geometry is the concept of a unit vector. A unit vector is a vector with a length or
magnitude of one. The unit vectors are different for different coordinates. In Cartesian coordinates the directions are x and y
usually denoted x ^ and y
^ . With the triangle above the letters referred to as a “hat”. The unit vectors in Cartesian coordinates
describe a circle known as the “unit circle” which has radius one. This can be seen by taking all the possible vectors of length one
at all the possible angles in this coordinate system and placing them on the coordinates. If you were to draw a line around
connecting all the heads of all the vectors together, you would get a circle of radius one.
learning objectives
Examine the applications of vectors in analyzing physical quantities
Use of Vectors
Vectors can be used to represent physical quantities. Most commonly in physics, vectors are used to represent displacement,
velocity, and acceleration. Vectors are a combination of magnitude and direction, and are drawn as arrows. The length represents
the magnitude and the direction of that quantity is the direction in which the vector is pointing. Because vectors are constructed this
way, it is helpful to analyze physical quantities (with both size and direction) as vectors.
Applications
In physics, vectors are useful because they can visually represent position, displacement, velocity and acceleration. When drawing
vectors, you often do not have enough space to draw them to the scale they are representing, so it is important to denote somewhere
3.2.8 https://phys.libretexts.org/@go/page/14446
what scale they are being drawn at. For example, when drawing a vector that represents a magnitude of 100, one may draw a line
that is 5 units long at a scale of . When the inverse of the scale is multiplied by the drawn magnitude, it should equal the actual
1
20
magnitude.
Velocity
Velocity is also defined in terms of a magnitude and direction. To say that something is gaining or losing velocity one must also say
how much and in what direction. For example, an airplane flying at 200 km
h
to the northeast can be represented by an vector
pointing in the northeast direction with a magnitude of 200 km
h
. In drawing the vector, the magnitude is only important as a way to
compare two vectors of the same units. So, if there were another airplane flying 100 km
h
to the southwest, the vector arrow should
be half as long and pointing in the direction of southwest.
Acceleration
Acceleration, being the time rate of change of velocity, is composed of a magnitude and a direction, and is drawn with the same
concept as a velocity vector. A value for acceleration would not be helpful in physics if the magnitude and direction of this
acceleration was unknown, which is why these vectors are important. In a free body diagram, for example, of an object falling, it
would be helpful to use an acceleration vector near the object to denote its acceleration towards the ground. If gravity is the only
force acting on the object, this vector would be pointing downward with a magnitude of 9.81 of 32.2 .
m
s2
ft
s2
Vector Diagram: Here is a man walking up a hill. His direction of travel is defined by the angle theta relative to the vertical axis
and by the length of the arrow going up the hill. He is also being accelerated downward by gravity.
Key Points
Vectors can be broken down into two components: magnitude and direction.
By taking the vector to be analyzed as the hypotenuse, the horizontal and vertical components can be found by completing a
right triangle. The bottom edge of the triangle is the horizontal component and the side opposite the angle is the vertical
component.
The angle that the vector makes with the horizontal can be used to calculate the length of the two components.
Scalars are physical quantities represented by a single number and no direction.
Vectors are physical quantities that require both magnitude and direction.
Examples of scalars include height, mass, area, and volume. Examples of vectors include displacement, velocity, and
acceleration.
To add vectors, lay the first one on a set of axes with its tail at the origin. Place the next vector with its tail at the previous
vector’s head. When there are no more vectors, draw a straight line from the origin to the head of the last vector. This line is the
sum of the vectors.
To subtract vectors, proceed as if adding the two vectors, but flip the vector to be subtracted across the axes and then join it tail
to head as if adding.
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Adding or subtracting any number of vectors yields a resultant vector.
Vectors can be decomposed into horizontal and vertical components.
Once the vectors are decomposed into components, the components can be added.
Adding the respective components of two vectors yields a vector which is the sum of the two vectors.
A vector is a quantity with both magnitude and direction.
A scalar is a quantity with only magnitude.
Multiplying a vector by a scalar is equivalent to multiplying the vector’s magnitude by the scalar. The vector lengthens or
shrinks but does not change direction.
A unit vector is a vector of magnitude ( length ) 1.
A scalar is a physical quantity that can be represented by a single number. Unlike vectors, scalars do not have direction.
Multiplying a vector by a scalar is the same as multiplying the vector’s magnitude by the number represented by the scalar.
Vectors are arrows consisting of a magnitude and a direction. They are used in physics to represent physical quantities that also
have both magnitude and direction.
Displacement is a physics term meaning the distance of an object from a reference point. Since the displacement contains two
pieces of information: the distance from the reference point and the direction away from the point, it is well represented by a
vector.
Velocity is defined as the rate of change in time of the displacement. To know the velocity of an object one must know both
how fast the displacement is changing and in what direction. Therefore it is also well represented by a vector.
Acceleration, being the rate of change of velocity also requires both a magnitude and a direction relative to some coordinates.
When drawing vectors, you often do not have enough space to draw them to the scale they are representing, so it is important to
denote somewhere what scale they are being drawn at.
Key Terms
coordinates: Numbers indicating a position with respect to some axis. Ex: x and y coordinates indicate position relative to xx
and yy axes.
axis: An imaginary line around which an object spins or is symmetrically arranged.
magnitude: A number assigned to a vector indicating its length.
Coordinate axes: A set of perpendicular lines which define coordinates relative to an origin. Example: x and y coordinate axes
define horizontal and vertical position.
origin: The center of a coordinate axis, defined as being the coordinate 0 in all axes.
Component: A part of a vector. For example, horizontal and vertical components.
vector: A directed quantity, one with both magnitude and direction; the between two points.
magnitude: A number assigned to a vector indicating its length.
scalar: A quantity that has magnitude but not direction; compare vector.
unit vector: A vector of magnitude 1.
velocity: The rate of change of displacement with respect to change in time.
displacement: The length and direction of a straight line between two objects.
acceleration: the rate at which the velocity of a body changes with time
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3.2.13 https://phys.libretexts.org/@go/page/14446
3.3: Projectile Motion
Basic Equations and Parabolic Path
Projectile motion is a form of motion where an object moves in parabolic path; the path that the object follows is called its
trajectory.
learning objectives
Assess the effect of angle and velocity on the trajectory of the projectile; derive maximum height using displacement
Projectile Motion
Projectile motion is a form of motion where an object moves in a bilaterally symmetrical, parabolic path. The path that the object
follows is called its trajectory. Projectile motion only occurs when there is one force applied at the beginning on the trajectory, after
which the only interference is from gravity. In a previous atom we discussed what the various components of an object in projectile
motion are. In this atom we will discuss the basic equations that go along with them in the special case in which the projectile
initial positions are null (i.e. x = 0 and y = 0 ).
0 0
Initial Velocity
The initial velocity can be expressed as x components and y components:
ux = u ⋅ cos θ (3.3.1)
uy = u ⋅ sin θ (3.3.2)
In this equation, u stands for initial velocity magnitude and θ refers to projectile angle.
Time of Flight
The time of flight of a projectile motion is the time from when the object is projected to the time it reaches the surface. As we
discussed previously, T depends on the initial velocity magnitude and the angle of the projectile:
2 ⋅ uy
T= (3.3.3)
g
2 ⋅ u ⋅ sin θ
T= (3.3.4)
g
Acceleration
In projectile motion, there is no acceleration in the horizontal direction. The acceleration, a, in the vertical direction is just due to
gravity, also known as free fall:
ax = 0 (3.3.5)
ay = −g (3.3.6)
Velocity
The horizontal velocity remains constant, but the vertical velocity varies linearly, because the acceleration is constant. At any time,
t , the velocity is:
ux = u ⋅ cos θ (3.3.7)
uy = u ⋅ sin θ − g ⋅ t (3.3.8)
Displacement
At time, t, the displacement components are:
3.3.1 https://phys.libretexts.org/@go/page/14447
x = u ⋅ t ⋅ cos θ (3.3.10)
1 2
y = u ⋅ t ⋅ sin θ − gt (3.3.11)
2
−− −−−−
The equation for the magnitude of the displacement is Δr = √x 2
+y
2
.
Parabolic Trajectory
We can use the displacement equations in the x and y direction to obtain an equation for the parabolic form of a projectile motion:
g 2
y = tan θ ⋅ x − ⋅x (3.3.12)
2 2
2⋅u ⋅ cos θ
Maximum Height
The maximum height is reached when v y =0 . Using this we can rearrange the velocity equation to find the time it will take for the
object to reach maximum height
u ⋅ sin θ
th = (3.3.13)
g
where t stands for the time it takes to reach maximum height. From the displacement equation we can find the maximum height
h
2 2
u ⋅ sin θ
h= (3.3.14)
2⋅g
Range
The range of the motion is fixed by the condition y =0 . Using this we can rearrange the parabolic motion equation to find the
range of the motion:
2
u ⋅ sin 2θ
R = . (3.3.15)
g
3.3.2 https://phys.libretexts.org/@go/page/14447
Projectiles at an Angle - Very Clear Explanation
Projectiles at an Angle: This video gives a clear and simple explanation of how to solve a problem on Projectiles Launched at an
Angle. I try to go step by step through this difficult problem to layout how to solve it in a super clear way. 2D kinematic problems
take time to solve, take notes on the order of how I solved it. Best wishes. Tune into my other videos for more help. Peace.
Solving Problems
In projectile motion, an object moves in parabolic path; the path the object follows is called its trajectory.
learning objectives
Identify which components are essential in determining projectile motion of an object
We have previously discussed projectile motion and its key components and basic equations. Using that information, we can solve
many problems involving projectile motion. Before we do this, let’s review some of the key factors that will go into this problem-
solving.
3.3.3 https://phys.libretexts.org/@go/page/14447
Initial velocity, u
Time of flight, T
Acceleration, a
Horizontal velocity, v x
Vertical velocity, v
y
Displacement, d
Maximum height, H
Range, R
How To Solve Any Projectile Motion Problem (The Toolbox Method): Introducing the “Toolbox” method of solving projectile
motion problems! Here we use kinematic equations and modify with initial conditions to generate a “toolbox” of equations with
which to solve a classic three-part projectile motion problem.
Now, let’s look at two examples of problems involving projectile motion.
Example 3.3.1:
Example 1
Let’s say you are given an object that needs to clear two posts of equal height separated by a specific distance. Refer to for this
–
example. The projectile is thrown at 25√2 m/s at an angle of 45°. If the object is to clear both posts, each with a height of
30m, find the minimum: (a) position of the launch on the ground in relation to the posts and (b) the separation between the
posts. For simplicity’s sake, use a gravity constant of 10. Problems of any type in physics are much easier to solve if you list
the things that you know (the “givens”).
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Diagram for Example 1: Use this figure as a reference to solve example 1. The problem is to make sure the object is able to
clear both posts.
Solution: The first thing we need to do is figure out at what time tt the object reaches the specified height. Since the motion is
in a parabolic shape, this will occur twice: once when traveling upward, and again when the object is traveling downward. For
this we can use the equation of displacement in the vertical direction, y − y : 0
1
2
y − y0 = (vy ⋅ t) − ( ⋅g⋅t ) (3.3.16)
2
Therefore:
1
2
30m = 25 ⋅ t − ⋅ 10 ⋅ t (3.3.18)
2
We can use the quadratic equation to find that the roots of this equation are 2s and 3s. This means that the projectile will reach
30m after 2s, on its way up, and after 3s, on its way down.
Example 2
An object is launched from the base of an incline, which is at an angle of 30°. If the launch angle is 60° from the horizontal and
the launch speed is 10 m/s, what is the total flight time? The following information is given: u = 10 ; θ = 60°; g = 10 .
m
s
m
s
2
Diagram for Example 2: When dealing with an object in projectile motion on an incline, we first need to use the given
information to reorient the coordinate system in order to have the object launch and fall on the same surface.
Solution: In order to account for the incline angle, we have to reorient the coordinate system so that the points of projection
and return are on the same level. The angle of projection with respect to the x direction is θ − α , and the acceleration in the y
direction is g ⋅ cos α . We replace θ with θ − α and g with g ⋅ cos α :
2 ⋅ u ⋅ sin(θ) 2 ⋅ u ⋅ sin(θ − α) 2 ⋅ 10 ⋅ sin(60 − 30) 20 ⋅ sin(30)
T = = = = (3.3.19)
g g ⋅ cos(α) 10 ⋅ cos(30) 10 ⋅ cos(30)
2
T = s (3.3.20)
–
√3
3.3.5 https://phys.libretexts.org/@go/page/14447
Zero Launch Angle
−
−− −
−−
An object launched horizontally at a height H travels a range v 0√
2H
g
during a time of flight T = √ 2H
g
.
learning objectives
Explain the relationship between the range and the time of flight
Projectile motion is a form of motion where an object moves in a parabolic path. The path followed by the object is called its
trajectory. Projectile motion occurs when a force is applied at the beginning of the trajectory for the launch (after this the projectile
is subject only to the gravity).
One of the key components of the projectile motion, and the trajectory it follows, is the initial launch angle. The angle at which the
object is launched dictates the range, height, and time of flight the object will experience while in projectile motion. shows
different paths for the same object being launched at the same initial velocity and different launch angles. As illustrated by the
figure, the larger the initial launch angle and maximum height, the longer the flight time of the object.
Projectile Trajectories: The launch angle determines the range and maximum height that an object will experience after being
launched.This image shows that path of the same object being launched at the same speed but different angles.
We have previously discussed the effects of different launch angles on range, height, and time of flight. However, what happens if
there is no angle, and the object is just launched horizontally? It makes sense that the object should be launched at a certain height (
H ), otherwise it wouldn’t travel very far before hitting the ground. Let’s examine how an object launched horizontally at a height H
α
X
g
Projectile motion: Projectile moving following a parabola.Initial launch angle is αα, and the velocity is v . 0
Duration of Flight
There is no vertical component in the initial velocity (v ) because the object is launched horizontally. Since the object travels
0
distance H in the vertical direction before it hits the ground, we can use the kinematic equation for the vertical motion:
3.3.6 https://phys.libretexts.org/@go/page/14447
1 2
(y − y0 ) = −H = 0 ⋅ T − gT (3.3.21)
2
Here, T is the duration of the flight before the object its the ground. Therefore:
−−−
2H
T =√ (3.3.22)
g
Range
In the horizontal direction, the object travels at a constant speed v0 during the flight. Therefore, the range R (in the horizontal
direction) is given as:
−−−
2H
R = v0 ⋅ T = v0 √ (3.3.23)
g
learning objectives
Choose the appropriate equation to find range, maximum height, and time of flight
Projectile motion is a form of motion where an object moves in a bilaterally symmetrical, parabolic path. The path that the object
follows is called its trajectory. Projectile motion only occurs when there is one force applied at the beginning of the trajectory, after
which the only interference is from gravity.
One of the key components of projectile motion and the trajectory that it follows is the initial launch angle. This angle can be
anywhere from 0 to 90 degrees. The angle at which the object is launched dictates the range, height, and time of flight it will
experience while in projectile motion. shows different paths for the same object launched at the same initial velocity at different
launch angles. As you can see from the figure, the larger the initial launch angle, the closer the object comes to maximum height
and the longer the flight time. The largest range will be experienced at a launch angle up to 45 degrees.
Launch Angle: The launch angle determines the range and maximum height that an object will experience after being launched.
This image shows that path of the same object being launched at the same velocity but different angles.
The range, maximum height, and time of flight can be found if you know the initial launch angle and velocity, using the following
equations:
2
v sin 2 θi
i
R = (3.3.24)
g
2 2
v sin θi
i
h = (3.3.25)
2g
2 vi sin θ
T = (3.3.26)
g
Where R – Range, h – maximum height, T – time of flight, vi – initial velocity, θi – initial launch angle, g – gravity.
3.3.7 https://phys.libretexts.org/@go/page/14447
Now that we understand how the launch angle plays a major role in many other components of the trajectory of an object in
projectile motion, we can apply that knowledge to making an object land where we want it. If there is a certain distance, d, that you
want your object to go and you know the initial velocity at which it will be launched, the initial launch angle required to get it that
distance is called the angle of reach. It can be found using the following equation:
1 gd
−1
θ = sin ( ) (3.3.27)
2
2 v
learning objectives
Construct a model of projectile motion by including time of flight, maximum height, and range
Time of Flight, T:
The time of flight of a projectile motion is exactly what it sounds like. It is the time from when the object is projected to the time it
reaches the surface. The time of flight depends on the initial velocity of the object and the angle of the projection, θθ. When the
point of projection and point of return are on the same horizontal plane, the net vertical displacement of the object is zero.
Symmetry:
All projectile motion happens in a bilaterally symmetrical path, as long as the point of projection and return occur along the same
horizontal surface. Bilateral symmetry means that the motion is symmetrical in the vertical plane. If you were to draw a straight
vertical line from the maximum height of the trajectory, it would mirror itself along this line.
Maximum Height, H:
The maximum height of a object in a projectile trajectory occurs when the vertical component of velocity, vyvy, equals zero. As the
projectile moves upwards it goes against gravity, and therefore the velocity begins to decelerate. Eventually the vertical velocity
will reach zero, and the projectile is accelerated downward under gravity immediately. Once the projectile reaches its maximum
height, it begins to accelerate downward. This is also the point where you would draw a vertical line of symmetry.
3.3.8 https://phys.libretexts.org/@go/page/14447
Range: The range of a projectile motion, as seen in this image, is independent of the forces of gravity.
Key Points
Objects that are projected from, and land on the same horizontal surface will have a vertically symmetrical path.
The time it takes from an object to be projected and land is called the time of flight. This depends on the initial velocity of the
projectile and the angle of projection.
When the projectile reaches a vertical velocity of zero, this is the maximum height of the projectile and then gravity will take
over and accelerate the object downward.
The horizontal displacement of the projectile is called the range of the projectile, and depends on the initial velocity of the
object.
When solving problems involving projectile motion, we must remember all the key components of the motion and the basic
equations that go along with them.
Using that information, we can solve many different types of problems as long as we can analyze the information we are given
and use the basic equations to figure it out.
To clear two posts of equal height, and to figure out what the distance between these posts is, we need to remember that the
trajectory is a parabolic shape and that there are two different times at which the object will reach the height of the posts.
When dealing with an object in projectile motion on an incline, we first need to use the given information to reorientate the
coordinate system in order to have the object launch and fall on the same surface.
For the zero launch angle, there is no vertical component in the initial velocity.
−
−−
The duration of the flight before the object hits the ground is given as T = √ 2H
g
.
In the horizontal direction, the object travels at a constant speed v0 during the flight. The range R (in the horizontal direction) is
−−−
2H
given as: R = v 0 ⋅ T = v0 √ .
g
If the same object is launched at the same initial velocity, the height and time of flight will increase proportionally to the initial
launch angle.
An object launched into projectile motion will have an initial launch angle anywhere from 0 to 90 degrees.
2
v sin 2 θi
i
The range of an object, given the initial launch angle and initial velocity is found with: R = .
g
2 2
v sin θi
i
The maximum height of an object, given the initial launch angle and initial velocity is found with: = .
2g
2 vi sin θ
The time of flight of an object, given the initial launch angle and initial velocity is found with: T = .
g
1 gd
The angle of reach is the angle the object must be launched at in order to achieve a specific distance: θ = sin
−1
(
2
) .
2 v
Objects that are projected from and land on the same horizontal surface will have a path symmetric about a vertical line through
a point at the maximum height of the projectile.
The time it takes from an object to be projected and land is called the time of flight. It depends on the initial velocity of the
projectile and the angle of projection.
The maximum height of the projectile is when the projectile reaches zero vertical velocity. From this point the vertical
component of the velocity vector will point downwards.
3.3.9 https://phys.libretexts.org/@go/page/14447
The horizontal displacement of the projectile is called the range of the projectile and depends on the initial velocity of the
object.
If an object is projected at the same initial speed, but two complementary angles of projection, the range of the projectile will be
the same.
Key Terms
trajectory: The path of a body as it travels through space.
symmetrical: Exhibiting symmetry; having harmonious or proportionate arrangement of parts; having corresponding parts or
relations.
reorientate: to orientate anew; to cause to face a different direction
gravity: Resultant force on Earth’s surface, of the attraction by the Earth’s masses, and the centrifugal pseudo-force caused by
the Earth’s rotation.
bilateral symmetry: the property of being symmetrical about a vertical plane
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How To Solve Any Projectile Motion Problem (The Toolbox Method). Located at: http://www.youtube.com/watch?
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How To Solve Any Projectile Motion Problem (The Toolbox Method). Located at: http://www.youtube.com/watch?
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3.4: Multiple Velocities
Addition of Velocities
Relative velocities can be found by adding the velocity of the observed object to the velocity of the frame of reference it was
measured in.
As learned in a previous atom, relative velocity is the velocity of an object as observed from a certain frame of reference.
demonstrates the concept of relative velocity. The girl is riding in a sled at 1.0 m/s, relative to an observer. When she throws the
snowball forward at a speed of 1.5 m/s, relative to the sled, the velocity of the snowball to the observer is the sum of the velocity of
the sled and the velocity of the snowball relative to the sled:
1.0m/s + 1.5m/s = 2.5m/s (3.4.1)
If the girl were to throw the snowball behind her at the same speed, the velocity of the ball relative to the observer would be:
1.0m/s − 1.5m/s = −0.5m/s (3.4.2)
The concept of relative velocity can also be demonstrated using the example of a boat in a river with a current. The boat is only
trying to move forward, but since the river is in motion, it carries the boat sideways while it moves forward. The person on the boat
is only observing the forward motion, while an observer on the shore will notice that the boat is moving sideways. In order to
calculate the velocity that the object is moving relative to earth, it is helpful to remember that velocity is a vector. In order to
analytically add these vectors, you need to remember the relationship between the magnitude and direction of the vector and its
components on the x and y axis of the coordinate system:
−−−−−−
Magnitude: v = √v 2
x + vy
2
vy
Direction: θ = tan ( )
−1
vx
x-component: v = v cos θ
x
y-component: v = v sin θ
y
Relative Velocity
Relative velocity is the velocity of an object B measured with respect to the velocity of another object A, denoted as v BA .
Relative velocity is the velocity of an object B, in the rest frame of another object A. This is denoted as vAB , where A is the
velocity in the rest frame of B.
Galileo observed the concept of relative velocity by using an example of a fly and a boat. He observed that while you are aboard
the boat, if you see a fly, you can measure its velocity, u. You can then go back on land and measure the velocity of the boat, v . Is
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the velocity of the fly, u, the actual velocity of the fly? No, because what you measured was the velocity of the fly relative to the
velocity of the boat. To obtain the velocity of the fly relative to the shore, s, you can use the vector sum as shown: s = u + v
Key Points
In order to find the velocity of an object B that is moving on object A that is observed by an observer that is not moving, add
the velocity of B and A together.
Velocity is a vector quantity, so the relationships between the magnitude, direction, x- axis component and y-axis component
are important.
These vector components can be added analytically or graphically.
In order to calculate the magnitude and direction, you must know the values of the axis components (either x and y , or x, y, and
z ) and to calculate the component values, you must know the magnitude and direction.
Relative velocity is the velocity of an object in motion being observed from a frame of reference that is either also in motion or
stationary.
If the frame of reference is moving in the same direction as the object being observed, it will appear as though the object is
moving slower than it actually is.
If the object being observed is on a moving surface, the velocity observed from that surface will be less than the velocity
observed from a stationary surface looking onto the moving surface.
In Galileo’s example, the observed velocity of the fly, uu, measured in reference to the velocity of the boat, v . In order to find
the velocity of the fly with respect to the shore, S, he had to add the velocity of the boat to the observed velocity of the fly:
s = u+v .
Key Terms
relative: Expressed in relation to another item, rather than in complete form.
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CHAPTER OVERVIEW
30: Nuclear Physics and Radioactivity is shared under a not declared license and was authored, remixed, and/or curated by LibreTexts.
1
30.1: The Nucleus
learning objectives
Explain relationship between nuclear radius, nuclear density, and nuclear size.
Nuclear size is defined by nuclear radius, also called rms charge radius. It can be measured by the scattering of electrons by the
nucleus and also inferred from the effects of finite nuclear size on electron energy levels as measured in atomic spectra.
The problem of defining a radius for the atomic nucleus is similar to the problem of atomic radius, in that neither atoms nor their
nuclei have definite boundaries. However, the nucleus can be modelled as a sphere of positive charge for the interpretation of
electron scattering experiments: because there is no definite boundary to the nucleus, the electrons “see” a range of cross-sections,
for which a mean can be taken. The qualification of “rms” (for “root mean square”) arises because it is the nuclear cross-section,
proportional to the square of the radius, which is determining for electron scattering.
The first estimate of a nuclear charge radius was made by Hans Geiger and Ernest Marsden in 1909, under the direction of Ernest
Rutherford at the Physical Laboratories of the University of Manchester, UK. The famous Rutherford gold foil experiment involved
the scattering of α-particles by gold foil, with some of the particles being scattered through angles of more than 90°, that is coming
back to the same side of the foil as the α-source, as shown in Figure 1. Rutherford was able to put an upper limit on the radius of
the gold nucleus of 34 femtometers (fm).
Later studies found an empirical relation between the charge radius and the mass number, A , for heavier nuclei
(A > 20) : R ≈ r ⋅ A where r is an empirical constant of 1.2–1.5 fm. This gives a charge radius for the gold nucleus
1/3
Nuclear Stability
The stability of an atom depends on the ratio and number of protons and neutrons, which may represent closed and filled quantum
shells.
learning objectives
Explain the relationship between the stability of an atom and its atomic structure.
The stability of an atom depends on the ratio of its protons to its neutrons, as well as on whether it contains a “magic number” of
neutrons or protons that would represent closed and filled quantum shells. These quantum shells correspond to energy levels within
the shell model of the nucleus. Filled shells, such as the filled shell of 50 protons in the element tin, confers unusual stability on the
nuclide. Of the 254 known stable nuclides, only four have both an odd number of protons and an odd number of neutrons:
hydrogen-2 (deuterium)
lithium-6
boron-10
nitrogen-14
Also, only four naturally occurring, radioactive odd-odd nuclides have a half-life greater than a billion years:
potassium-40
vanadium-50
lanthanum-138
tantalum-180m
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Alpha Decay: Alpha decay is one type of radioactive decay. An atomic nucleus emits an alpha particle and thereby transforms
(“decays”) into an atom with a mass number smaller by four and an atomic number smaller by two. Many other types of decay are
possible.
Most odd-odd nuclei are highly unstable with respect to beta decay because the decay products are even-even and therefore more
strongly bound, due to nuclear pairing effects.
An atom with an unstable nucleus, called a radionuclide, is characterized by excess energy available either for a newly created
radiation particle within the nucleus or via internal conversion. During this process, the radionuclide is said to undergo radioactive
decay. Radioactive decay results in the emission of gamma rays and/or subatomic particles such as alpha or beta particles, as shown
in. These emissions constitute ionizing radiation. Radionuclides occur naturally but can also be produced artificially.
All elements form a number of radionuclides, although the half-lives of many are so short that they are not observed in nature.
Even the lightest element, hydrogen, has a well-known radioisotope: tritium. The heaviest elements (heavier than bismuth) exist
only as radionuclides. For every chemical element, many radioisotopes that do not occur in nature (due to short half-lives or the
lack of a natural production source) have been produced artificially.
learning objectives
Explain how nuclear force varies with distance.
The nuclear force is the force between two or more component parts of an atomic nuclei. The component parts are neutrons and
protons, which collectively are called nucleons. Nuclear force is responsible for the binding of protons and neutrons into atomic
nuclei.
Drawing of Atomic Nucleus: A model of the atomic nucleus showing it as a compact bundle of the two types of nucleons: protons
(red) and neutrons (blue).
To disassemble a nucleus into unbound protons and neutrons would require working against the nuclear force. Conversely, energy
is released when a nucleus is created from free nucleons or other nuclei—known as the nuclear binding energy. The binding energy
of nuclei is always a positive number, since all nuclei require net energy to separate into individual protons and neutrons. Because
of mass-energy equivalence (i.e., Einstein’s famous formula E = mc ), releasing this energy causes the mass of the nucleus to be
2
lower than the total mass of the individual nucleons (leading to “mass deficit”). Binding energy is the energy used in nuclear power
plants and nuclear weapons.
The nuclear force is powerfully attractive between nucleons at distances of about 1 femtometer (fm) between their centers, but
rapidly decreases to relative insignificance at distances beyond about 2.5 fm. At very short distances (less than 0.7 fm) it becomes
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repulsive; it is responsible for the physical size of nuclei since the nucleons can come no closer than the force allows.
The nuclear force is now understood as a residual effect of an even more powerful “strong force” or strong interaction. It is the
attractive force that binds together particles known as quarks (to form the nucleons themselves). This more powerful force is
mediated by particles called gluons. Gluons hold quarks together with a force like that of an electric charge (but of far greater
power).
The nuclear forces arising between nucleons are now seen as analogous to the forces in chemistry between neutral atoms or
molecules (called London forces). Such forces between atoms are much weaker than the attractive electrical forces that hold
together the atoms themselves (i.e., that bind electrons to the nucleus), and their range between atoms is shorter because they arise
from a small separation of charges inside the neutral atom.
Similarly, even though nucleons are made of quarks in combinations which cancel most gluon forces (they are “color neutral”),
some combinations of quarks and gluons leak away from nucleons in the form of short-range nuclear force fields that extend from
one nucleon to another nucleon in close proximity. These nuclear forces are very weak compared to direct gluon forces (“color
forces” or “strong forces”) inside nucleons, and the nuclear forces extend over only a few nuclear diameters, falling exponentially
with distance. Nevertheless, they are strong enough to bind neutrons and protons over short distances, as well as overcome the
electrical repulsion between protons in the nucleus. Like London forces, nuclear forces also stop being attractive, and become
repulsive when nucleons are brought too close together.
Key Points
The first estimate of a nuclear charge radius was made by Hans Geiger and Ernest Marsden in 1909, under the direction of
Ernest Rutherford, in the gold foil experiment that involved the scattering of α-particles by gold foil, as shown in Figure 1.
An empirical relation exists between the charge radius and the mass number, A , for heavier nuclei (A > 20 ): where r is an
empirical constant of 1.2–1.5 fm.
The nuclear density for a typical nucleus can be approximately calculated from the size of the nucleus: n = 4
A
3
πR
3
Most odd-odd nuclei are highly unstable with respect to beta decay because the decay products are even-even and therefore
more strongly bound, due to nuclear pairing effects.
An atom with an unstable nucleus is characterized by excess energy available either for a newly created radiation particle within
the nucleus or via internal conversion.
All elements form a number of radionuclides, although the half-lives of many are so short that they are not observed in nature.
The nuclear force is powerfully attractive at distances of about 1 femtometer (fm), rapidly decreases to insignificance at
distances beyond about 2.5 fm, and becomes repulsive at very short distances less than 0.7 fm.
The nuclear force is a residual effect of the a strong interaction that binds together particles called quarks into nucleons.
The binding energy of nuclei is always a positive number while the mass of an atom ‘s nucleus is always less than the sum of
the individual masses of the constituent protons and neutrons when separated.
Key Terms
α-particle: two protons and two neutrons bound together into a particle identical to a helium nucleus
atomic spectra: emission or absorption lines formed when an electron makes a transition from one energy level of an atom to
another
nucleus: the massive, positively charged central part of an atom, made up of protons and neutrons
nuclide: A nuclide (from “nucleus”) is an atomic species characterized by the specific constitution of its nucleus — i.e., by its
number of protons (Z), its number of neutrons (N ), and its nuclear energy state.
radionuclide: A radionuclide is an atom with an unstable nucleus, characterized by excess energy available to be imparted
either to a newly created radiation particle within the nucleus or via internal conversion.
radioactive decay: any of several processes by which unstable nuclei emit subatomic particles and/or ionizing radiation and
disintegrate into one or more smaller nuclei
nucleus: the massive, positively charged central part of an atom, made up of protons and neutrons
quark: In the Standard Model, an elementary subatomic particle that forms matter. Quarks are never found alone in nature, but
combine to form hadrons, such as protons and neutrons.
gluon: A massless gauge boson that binds quarks together to form baryons, mesons and other hadrons; it is associated with the
strong nuclear force.
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30.2: Radioactivity
Learning objectives
Name major sources of terrestrial radiation.
Radioactive material is found throughout nature. Detectable amounts occur naturally in soil, rocks, water, air, and vegetation. From
these sources it can be inhaled and ingested into the body. In addition to this internal exposure, humans also receive external
exposure from radioactive materials that remain outside the body and from cosmic radiation from space. The worldwide average
natural dose to humans is about 2.4 millisieverts (mSv) per year. This is four times more than the worldwide average artificial
radiation exposure, which in the year 2008 amounted to about 0.6 mSv per year. In some wealthier countries, such as the US and
Japan, artificial exposure is, on average, greater than the natural exposure, due to greater access to medical imaging. In Europe, the
average natural background exposure by country ranges from under 2 mSv annually in the United Kingdom to more than 7 mSv
annually in Finland, as shown in.
Natural Radiation Atlas of Europe: Bar chart of average annual dosages from natural radiation sources for major European
countries
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Terrestrial Radiation
Terrestrial radiation only includes sources that remain external to the body. The major radionuclides of concern are potassium,
uranium, and thorium and their decay products. Some of these decay products, like radium and radon, are intensely radioactive but
occur in low concentrations. Most of these sources have been decreasing, due to radioactive decay since the formation of the earth,
because there is no significant source of replacement. Because of this, the present activity on Earth from uranium-238 is only half
as much as it originally was because of its 4.5-billion-year half-life. Potassium-40 (with a half-life of 1.25 billion years) is at about
eight percent of its original activity. However, the effects on humans of the actual diminishment (due to decay) of these isotopes is
minimal. This is because humans evolved too recently for the difference in activity over a fraction of a half-life to be significant.
Put another way, human history is so short in comparison to a half-life of a billion years that the activity of these long-lived
isotopes has been effectively constant throughout our time on this planet.
Many shorter-half-life and therefore more intensely radioactive isotopes have not decayed out of the terrestrial environment
because they are still being produced. Examples of these are radium-226 (a decay product of uranium-238) and radon-222 (a decay
product of radium-226).
Radiation Detection
A radiation detector is a device used to detect, track, or identify high-energy particles.
Learning objectives
Explain difference between major types of radiation detectors.
A radiation detector is a device used to detect, track, or identify high- energy particles, such as those produced by nuclear decay,
cosmic radiation, and reactions in a particle accelerator. Modern detectors are also used as calorimeters to measure the energy of
detected radiation. They may be also used to measure other attributes, such as momentum, spin, and charge of the particles.
Different types of radiation detectors exist; gaseous ionization detectors, semiconductor detectors, and scintillation detectors are the
most common.
Semiconductor Detectors
A semiconductor detector uses a semiconductor (usually silicon or germanium) to detect traversing charged particles or the
absorption of photons. When these detectors’ sensitive structures are based on single diodes, they are called semiconductor diode
detectors. When they contain many diodes with different functions, the more general term “semiconductor detector” is used.
Semiconductor detectors have had various applications in recent decades, in particular in gamma and x-ray spectrometry and as
particle detectors.
30.2.2 https://phys.libretexts.org/@go/page/16243
Scintillation Detectors
A scintillation detector is created by coupling a scintillator — a material that exhibits luminescence when excited by ionizing
radiation — to an electronic light sensor, such as a photomultiplier tube (PMT) or a photodiode. PMTs absorb the light emitted by
the scintillator and re-emit it in the form of electrons via the photoelectric effect. The subsequent multiplication of those electrons
(sometimes called photo-electrons) results in an electrical pulse, which can then be analyzed. The pulse yields meaningful
information about the particle that originally struck the scintillator.
Scintillators are used by the American government, particularly Homeland Security, as radiation detectors. Scintillators can also be
used in neutron and high-energy particle physics experiments, new energy resource exploration, x-ray security, nuclear cameras,
computed tomography, and gas exploration. Other applications of scintillators include CT scanners and gamma cameras in medical
diagnostics, screens in computer monitors, and television sets.
Learning objectives
Describe importance of radioactive decay series for decay process.
Radioactive decay series, or decay chains, describe the radioactive decay of different discrete radioactive decay products as a
chained series of transformations. Most radioactive elements do not decay directly to a stable state; rather, they undergo a series of
decays until eventually a stable isotope is reached.
Radioactive Decay Series Diagram: This diagram provides examples of four decay series: thorium (in blue), radium (in red),
actinium (in green), and neptunium (in purple).
Decay stages are referred to by their relationship to previous or subsequent stages. A parent isotope is one that undergoes decay to
form a daughter isotope. The daughter isotope may be stable, or it may itself decay to form a daughter isotope of its own. The
daughter of a daughter isotope is sometimes called a granddaughter isotope.
The time it takes for a single parent atom to decay to an atom of its daughter isotope can vary widely, not only for different parent-
daughter chains, but also for identical pairings of parent and daughter isotopes. While the decay of a single atom occurs
spontaneously, the decay of an initial population of identical atoms over time, tt, follows a decaying exponential distribution,
e−te−t, where λλis called the decay constant. Because of this exponential nature, one of the properties of an isotope is its half-life,
the time by which half of an initial number of identical parent radioisotopes have decayed to their daughters. Half-lives have been
determined in laboratories for thousands of radioisotopes (radionuclides). These half-lives can range from nearly nonexistent spans
of time to as much as 10 years or more.
19
The intermediate stages often emit more radioactivity than the original radioisotope. When equilibrium is achieved, a
granddaughter isotope is present in proportion to its half-life. But, since its activity is inversely proportional to its half-life, any
nuclide in the decay chain finally contributes as much as the head of the chain. For example, natural uranium is not significantly
radioactive, but pitchblende, a uranium ore, is 13 times more radioactive because of the radium and other daughter isotopes it
contains. Not only are unstable radium isotopes significant radioactivity emitters, but as the next stage in the decay chain they also
generate radon, a heavy, inert, naturally occurring radioactive gas. Rock containing thorium and/or uranium (such as some granites)
emits radon gas, which can accumulate in enclosed places such as basements or underground mines. Radon exposure is considered
the leading cause of lung cancer in non-smokers.
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Alpha Decay
In alpha decay an atomic nucleus emits an alpha particle and transforms into an atom with smaller mass (by four) and atomic
number (by two).
Learning objectives
Describe the process, penetration power, and effects of alpha radiation
Alpha decay is a type of radioactive decay in which an atomic nucleus emits an alpha particle that consists of two protons and two
neutrons, as shown in. As the result of this process, the parent atom transforms (“decays”) into a new atom with a mass number
smaller by four and an atomic number smaller by two.
Alpha Decay: Alpha decay is one type of radioactive decay. An atomic nucleus emits an alpha particle and thereby transforms
(“decays”) into an atom with a mass number smaller by four and an atomic number smaller by two. Many other types of decay are
possible.
For example: 238U → 234Th + α
Because an alpha particle is the same as a helium-4 nucleus, which has mass number 4 and atomic number 2, this can also be
written as:
238
92U → 234
90Th + 4
2He
The alpha particle also has charge +2, but the charge is usually not written in nuclear equations, which describe nuclear reactions
without considering the electrons. This convention is not meant to imply that the nuclei necessarily occur in neutral atoms.
Alpha decay is by far the most common form of cluster decay, in which the parent atom ejects a defined daughter collection of
nucleons, leaving another defined product behind (in nuclear fission, a number of different pairs of daughters of approximately
equal size are formed). Alpha decay is the most common cluster decay because of the combined extremely high binding energy and
relatively small mass of the helium-4 product nucleus (the alpha particle).
Alpha decay typically occurs in the heaviest nuclides. In theory it can occur only in nuclei somewhat heavier than nickel (element
28), in which overall binding energy per nucleon is no longer a minimum and the nuclides are therefore unstable toward
spontaneous fission-type processes. The lightest known alpha emitters are the lightest isotopes (mass numbers 106-110) of
tellurium (element 52).
Alpha particles have a typical kinetic energy of 5 MeV (approximately 0.13 percent of their total energy, i.e., 110 TJ/kg) and a
speed of 15,000 km/s. This corresponds to a speed of around 0.05 c. There is surprisingly small variation in this energy, due to the
heavy dependence of the half-life of this process on the energy produced.
Because of their relatively large mass, +2 electric charge, and relatively low velocity, alpha particles are very likely to interact with
other atoms and lose their energy, so their forward motion is effectively stopped within a few centimeters of air.
Most of the helium produced on Earth (approximately 99 percent of it) is the result of the alpha decay of underground deposits of
minerals containing uranium or thorium. The helium is brought to the surface as a byproduct of natural gas production.
30.2.4 https://phys.libretexts.org/@go/page/16243
Beta Decay
Beta decay is a type of radioactive decay in which a beta particle (an electron or a positron) is emitted from an atomic nucleus.
Learning objectives
Explain difference between beta minus and beta plus decays.
Beta decay is a type of radioactive decay in which a beta particle (an electron or a positron) is emitted from an atomic nucleus, as
shown in. Beta decay is a process that allows the atom to obtain the optimal ratio of protons and neutrons.
Beta Decay: β decay in an atomic nucleus (the accompanying antineutrino is omitted). The inset shows beta decay of a free
neutron
There are two types of beta decay. Beta minus (β) leads to an electron emission (e−); beta plus (β+) leads to a positron emission
(e+). In electron emission an electron antineutrino is also emitted, while positron emission is accompanied by an electron neutrino.
Beta decay is mediated by the weak force.
Emitted beta particles have a continuous kinetic energy spectrum, ranging from 0 to the maximal available energy (Q), that depends
on the parent and daughter nuclear states that participate in the decay. The continuous energy spectra of beta particles occur
because Q is shared between a beta particle and a neutrino. A typical Q is around 1 MeV, but it can range from a few keV to a
several tens of MeV. Since the rest mass energy of the electron is 511 keV, the most energetic beta particles are ultrarelativistic,
with speeds very close to the speed of light.
Since the proton and neutron are part of an atomic nucleus, beta decay processes result in transmutation of one chemical element
into another. For example:
−
137Cs → 137Ba + e (30.2.1)
+
11Na → 10Ne + e (30.2.2)
Beta decay does not change the number of nucleons, A, in the nucleus; it changes only its charge, Z. Therefore the set of all
nuclides with the same A can be introduced; these isobaric nuclides may turn into each other via beta decay.
A beta-stable nucleus may undergo other kinds of radioactive decay (for example, alpha decay). In nature, most isotopes are beta-
stable, but there exist a few exceptions with half-lives so long that they have not had enough time to decay since the moment of
their nucleosynthesis. One example is the odd-proton odd-neutron nuclide 40 K, which undergoes both types of beta decay with a
half-life of 1.277 ·109 years.
30.2.5 https://phys.libretexts.org/@go/page/16243
Beta Decay 1/2
Beta Decay 1/2: In this video I introduce Beta decay and discuss it from an basic level to a perhaps second or third year University
level.
30.2.6 https://phys.libretexts.org/@go/page/16243
Beta Decay 2/2
Beta Decay 2/2: In this video I introduce Beta decay and discuss it from an basic level to a perhaps second or third year University
level.
Gamma Decay
Gamma decay is a process of emission of gamma rays that accompanies other forms of radioactive decay, such as alpha and beta
decay.
Learning objectives
Explain relationship between gamma decay and other forms of nuclear decay.
Gamma radiation, also known as gamma rays and denoted as γγ, is electromagnetic radiation of high frequency and therefore high
energy. Gamma rays typically have frequencies above 10 exahertz (>1019>1019 Hz) and therefore have energies above 100 keV
and wavelengths less than 10 picometers (less than the diameter of an atom). However, this is not a strict definition; rather, it is a
rule-of-thumb description for natural processes. Gamma rays from radioactive decay are defined as gamma rays no matter what
their energy, so there is no lower limit to gamma energy derived from radioactive decay. Gamma decay commonly produces
energies of a few hundred keV and usually less than 10 MeV.
30.2.7 https://phys.libretexts.org/@go/page/16243
Cobalt-60 Decay Scheme: Path of decay of Co-60 to Ni-60. Excited levels for Ni-60 that drop to ground state via emission of
gamma rays are indicated
Gamma decay accompanies other forms of decay, such as alpha and beta decay; gamma rays are produced after the other types of
decay occur. When a nucleus emits an α or β particle, the daughter nucleus is usually left in an excited state. It can then move to a
lower energy state by emitting a gamma ray, in much the same way that an atomic electron can jump to a lower energy state by
emitting a photon. For example, cobalt-60 decays to excited nickel-60 by beta decay through emission of an electron of 0.31 MeV.
Next, the excited nickel-60 drops down to the ground state by emitting two gamma rays in succession (1.17 MeV, then 1.33 MeV),
as shown in. Emission of a gamma ray from an excited nuclear state typically requires only 10−1210−12 seconds: it is nearly
instantaneous. Gamma decay from excited states may also follow nuclear reactions such as neutron capture, nuclear fission, or
nuclear fusion.
In certain cases, the excited nuclear state following the emission of a beta particle may be more stable than average; in these cases
it is termed a metastable excited state if its decay is 100 to 1000 times longer than the average 10−1210−12 seconds. Such nuclei
have half-lives that are easily measurable; these are termed nuclear isomers. Some nuclear isomers are able to stay in their excited
state for minutes, hours, or days, or occasionally far longer, before emitting a gamma ray. This phenomenon is called isomeric
transition. The process of isomeric transition is therefore similar to any gamma emission; it differs only in that it involves the
intermediate metastable excited states of the nuclei.
Learning objectives
Identify the age of materials that can be approximately determined using radiocarbon dating
Radiocarbon dating (usually referred to simply as carbon-14 dating) is a radiometric dating method. It uses the naturally occurring
radioisotope carbon-14 (14C) to estimate the age of carbon-bearing materials up to about 58,000 to 62,000 years old.
Carbon has two stable, nonradioactive isotopes: carbon-12 (12C) and carbon-13 (13C). There are also trace amounts of the unstable
radioisotope carbon-14 (14C) on Earth. Carbon-14 has a relatively short half-life of 5,730 years, meaning that the fraction of
carbon-14 in a sample is halved over the course of 5,730 years due to radioactive decay to nitrogen-14. The carbon-14 isotope
would vanish from Earth’s atmosphere in less than a million years were it not for the constant influx of cosmic rays interacting with
molecules of nitrogen (N2) and single nitrogen atoms (N) in the stratosphere. Both processes of formation and decay of carbon-14
are shown in.
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Formation and Decay of Carbon-14: Diagram of the formation of carbon-14 (1), the decay of carbon-14 (2), and equations
describing the carbon-12:carbon-14 ratio in living and dead organisms
When plants fix atmospheric carbon dioxide (CO2) into organic compounds during photosynthesis, the resulting fraction of the
isotope 14C in the plant tissue will match the fraction of the isotope in the atmosphere. After plants die or are consumed by other
organisms, the incorporation of all carbon isotopes, including 14C, stops. Thereafter, the concentration (fraction) of 14C declines at
a fixed exponential rate due to the radioactive decay of 14C. (An equation describing this process is shown in. ) Comparing the
remaining 14C fraction of a sample to that expected from atmospheric 14C allows us to estimate the age of the sample.
Raw (i.e., uncalibrated) radiocarbon ages are usually reported in radiocarbon years “Before Present” (BP), with “present” defined
as CE 1950. Such raw ages can be calibrated to give calendar dates. One of the most frequent uses of radiocarbon dating is to
estimate the age of organic remains from archaeological sites.
The technique of radiocarbon dating was developed by Willard Libby and his colleagues at the University of Chicago in 1949.
Emilio Segrè asserted in his autobiography that Enrico Fermi suggested the concept to Libby at a seminar in Chicago that year.
Libby estimated that the steady-state radioactivity concentration of exchangeable carbon-14 would be about 14 disintegrations per
minute (dpm) per gram. In 1960, Libby was awarded the Nobel Prize in chemistry for this work. He demonstrated the accuracy of
radiocarbon dating by accurately estimating the age of wood from a series of samples for which the age was known, including an
ancient Egyptian royal barge dating from 1850 BCE.
30.2.9 https://phys.libretexts.org/@go/page/16243
Half-life (2 of 16) An Explanation
Half-life: Describes radioactive half life and how to do some simple calculations using half life.
Learning objectives
Explain what is a half-life of a radionuclide.
The half-life of a radionuclide is the time taken for half of the radionuclide’s atoms to decay. Taking λλ to be the decay rate
(number of disintegrations per unit time), and ττ the average lifetime of an atom before it decays, we have:
−λt −t/τ
N(t) = N0 e = N0 e (30.2.3)
The half-life is related to the decay constant by substituting the condition N = N o /2 and solving for t = t
1/2 :
A half-life must not be thought of as the time required for exactly half of the atoms to decay.
30.2.10 https://phys.libretexts.org/@go/page/16243
Radioactive decay simulation: A simulation of many identical atoms undergoing radioactive decay, starting with four atoms (left)
and 400 atoms (right). The number at the top indicates how many half-lives have elapsed
The following figure shows a simulation of many identical atoms undergoing radioactive decay. Note that after one half-life there
are not exactly one-half of the atoms remaining; there are only approximatelyone-half left because of the random variation in the
process. However, with more atoms (the boxes on the right), the overall decay is smoother and less random-looking than with
fewer atoms (the boxes on the left), in accordance with the law of large numbers.
The relationship between the half-life and the decay constant shows that highly radioactive substances are quickly spent while
those that radiate weakly endure longer. Half-lives of known radionuclides vary widely, from more than 1019 years, such as for the
very nearly stable nuclide 209 Bi, to 10−23 seconds for highly unstable ones.
The factor of ln(2) in the above equations results from the fact that the concept of “half-life” is merely a way of selecting a different
base other than the natural base e for the lifetime expression. The time constant τis the e-1-life, the time until only 1/e remains —
about 36.8 percent, rather than the 50 percent in the half-life of a radionuclide. Therefore, τis longer than t1/2. The following
equation can be shown to be valid:
−t/τ −t/t1/2
N(t) = N0 e = N0 2 (30.2.5)
Since radioactive decay is exponential with a constant probability, each process could just as easily be described with a different
constant time period that (for example) gave its 1/3-life (how long until only 1/3 is left), or its 1/10-life (how long until only 1/10 is
left), and so on. Therefore, the choice of τand t1/2 for marker-times is only for convenience and for the sake of uploading
convention. These marker-times reflect a fundamental principle only in that they show that the same proportion of a given
radioactive substance will decay over any time period you choose.
Mathematically, the nth life for the above situation would be found by the same process shown above — by setting N = N 0 /n and
substituting into the decay solution, to obtain:
ln n
t1/n = = τ ln n (30.2.6)
λ
30.2.11 https://phys.libretexts.org/@go/page/16243
Joy physics probs 27- Half-life
Half-life: Part of a series of videos on physics problem-solving. The problems are taken from “The Joy of Physics. ” This one deals
with radioactive half-life. The viewer is urged to pause the video at the problem statement and work the problem before watching
the rest of the video.
Key Points
The biggest source of natural background radiation is airborne radon, a radioactive gas that emanates from the ground.
The earth is constantly bombarded by radiation from outer space that consists of positively charged ions ranging from protons
to iron and larger nuclei from sources outside of our solar system.
Terrestrial radiation includes sources that remain external to the body. The major radionuclides of concern are potassium,
uranium, and thorium and their decay products.
Gaseous ionization detectors use the ionizing effect of radiation upon gas-filled sensors.
A semiconductor detector uses a semiconductor (usually silicon or germanium) to detect traversing charged particles or the
absorption of photons.
A scintillation detector is created by coupling a scintillator to an electronic light sensor.
Most radioactive elements do not decay directly to a stable state; rather, they undergo a series of decays until eventually a stable
isotope is reached.
Half-lives of radioisotopes range from nearly nonexistent spans of time to as much as 1019 years or more.
The intermediate stages of radioactive decay series often emit more radioactivity than the original radioisotope.
An alpha particle is the same as a helium-4 nucleus, which has mass number 4 and atomic number 2.
Because of their relatively large mass, +2 electric charge, and relatively low velocity, alpha particles are very likely to interact
with other atoms and lose their energy, so their forward motion is effectively stopped within a few centimeters of air.
30.2.12 https://phys.libretexts.org/@go/page/16243
Most of the helium produced on Earth (approximately 99 percent of it) is the result of the alpha decay of underground deposits
of minerals containing uranium or thorium.
There are two types of beta decay: beta minus, which leads to an electron emission, and beta plus, which leads to a positron
emission.
Beta decay allows the atom to obtain the optimal ratio of protons and neutrons.
Beta decay processes transmute one chemical element into another.
Gamma decay accompanies other forms of decay, such as alpha and beta decay; gamma rays are produced after the other types
of decay occur.
Although emission of gamma ray is a nearly instantaneous process, it can involve intermediate metastable excited states of the
nuclei.
Gamma rays are generally the most energetic form of electromagnetic radiation.
Carbon-14 dating can be used to estimate the age of carbon-bearing materials up to about 58,000 to 62,000 years old.
The carbon-14 isotope would vanish from Earth’s atmosphere in less than a million years were it not for the constant influx of
cosmic rays interacting with atmospheric nitrogen.
One of the most frequent uses of radiocarbon dating is to estimate the age of organic remains from archaeological sites.
The half-life is related to the decay constant as follows: t = ln 2/λ .
1/2
The relationship between the half-life and the decay constant shows that highly radioactive substances are quickly spent while
those that radiate weakly endure longer.
Half-lives of known radionuclides vary widely, from more than 1019 years, such as for the very nearly stable nuclide 209 Bi, to
10-23 seconds for highly unstable ones.
Key Terms
radionuclide: A radionuclide is an atom with an unstable nucleus, characterized by excess energy available to be imparted
either to a newly created radiation particle within the nucleus or via internal conversion.
radon: a radioactive chemical element (symbol Rn, formerly Ro) with atomic number 86; one of the noble gases
sievert: in the International System of Units, the derived unit of radiation dose; the dose received in one hour at a distance of 1
cm from a point source of 1 mg of radium in a 0.5 mm thick platinum enclosure; symbol: Sv
scintillator: any substance that glows under the action of photons or other high-energy particles
diode: an electronic device that allows current to flow in one direction only; a valve
semiconductor: A substance with electrical properties intermediate between a good conductor and a good insulator.
half-life: the time required for half of the nuclei in a sample of a specific isotope to undergo radioactive decay
radioisotope: a radioactive isotope of an element
decay: to change by undergoing fission, by emitting radiation, or by capturing or losing one or more electrons
alpha particle: A positively charged nucleus of a helium-4 atom (consisting of two protons and two neutrons), emitted as a
consequence of radioactivity; α-particle.
radioactive decay: any of several processes by which unstable nuclei emit subatomic particles and/or ionizing radiation and
disintegrate into one or more smaller nuclei
beta decay: a nuclear reaction in which a beta particle (electron or positron) is emitted
positron: The antimatter equivalent of an electron, having the same mass but a positive charge.
transmutation: the transformation of one element into another by a nuclear reaction
electromagnetic radiation: radiation (quantized as photons) consisting of oscillating electric and magnetic fields oriented
perpendicularly to each other, moving through space
gamma ray: A very high frequency (and therefore very high energy) electromagnetic radiation emitted as a consequence of
radioactivity.
radiometric dating: Radiometric dating is a technique used to date objects based on a comparison between the observed
abundance of a naturally occurring radioactive isotope and its decay products using known decay rates.
carbon-14: carbon-14 is a radioactive isotope of carbon with a nucleus containing 6 protons and 8 neutrons.
radionuclide: A radionuclide is an atom with an unstable nucleus, characterized by excess energy available to be imparted
either to a newly created radiation particle within the nucleus or via internal conversion.
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30.3: Quantum Tunneling and Conservation Laws
learning objectives
Identify factors that affect the tunneling probability
Imagine throwing a ball at a wall and having it disappear the instant before it makes contact and appear on the other side. The wall
remains intact; the ball did not break through it. Believe it or not, there is a finite (if extremely small) probability that this event
would occur. This phenomenon is called quantum tunneling.
While the possibility of tunneling is essentially ignorable at macroscopic levels, it occurs regularly on the nanoscale level.
Consider, for example, a p-orbital in an atom. Between the two lobes there is a nodal plane. By definition there is precisely 0
probability of finding an electron anywhere along that plane, and because the plane extends infinitely it is impossible for an
electron to go around it. Yet, electrons commonly cross from one lobe to the other via quantum tunneling. They never exist in the
nodal area (this is forbidden); instead they travel through imaginary space.
P-Orbital: The red and blue lobes represent the volume in which there is a 90 percent probability of finding an electron at any
given time if the orbital is occupied.
Imaginary space is not real, but it is explicitly referenced in the time-dependent Schrödinger equation, which has a component of ii
(the square root of −1−1, an imaginary number):
∂
^
iℏ Ψ = HΨ (30.3.1)
∂t
And because all matter has a wave component (see the topic of wave-particle duality), all matter can in theory exist in imaginary
space. But what accounts for the difference in probability of an electron tunneling over a nodal plane and a ball tunneling through a
brick wall? The answer is a combination of the tunneling object’s mass (mm) and energy (EE) and the energy height (U0U0) of the
barrier through which it must travel to get to the other side.
When it reaches a barrier it cannot overcome, a particle’s wave function changes from sinusoidal to exponentially diminishing in
form. The solution for the Schrödinger equation in such a medium is:
−αx
Ψ = Ae (30.3.2)
where:
−−−−−−−−−−−
2m (U0 − E)
α =√ (30.3.3)
2
h
Therefore, the probability of an object tunneling through a barrier decreases with the object’s increasing mass and with the
increasing gap between the energy of the object and the energy of the barrier. And although the wave function never quite reaches 0
(as can be determined from the e−xe−x functionality), this explains how tunneling is frequent on nanoscale but negligible at the
macroscopic level.
30.3.1 https://phys.libretexts.org/@go/page/16244
Conservation of Nucleon Number and Other Laws
Through radioactive decay, nuclear fusion and nuclear fission, the number of nucleons (sum of protons and neutrons) is always
held constant.
learning objectives
Define the Law of Conservation of Nuclear Number
In physics and chemistry there are many conservation laws—among them, the Law of Conservation of Nucleon Number, which
states that the total number of nucleons (nuclear particles, specifically protons and neutrons) cannot change by any nuclear reaction.
Radioactive Decay
Consider the three modes of decay. In gamma decay, an excited nucleus releases gamma rays, but its proton (Z) and neutron (A-Z)
count remain the same:
A A
→ +γ (30.3.4)
ZN ZN
Nuclear fission of U-235: If U-235 is bombarded with a neutron (light blue small circe), the resulting U-236 produced is unstable
and undergoes fission. The resulting elements (shown here as Kr-92 and Ba-141) do not contain as many nucleons as U-236, with
the remaining three neutrons being released as high-energy particles, able to bombard another U-235 atom and maintain a chain
reaction.
In beta decay, a nucleus releases energy and either an electron or a positron. In the case of an electron being released, atomic mass
(A) remains the same as a neutron is converted into a proton, raising atomic number by 1:
A A −
→ N+e +ν (30.3.5)
ZN Z+1
In the case of a positron being released, atomic mass remains constant as a proton is converted to a neutron, lowering atomic
number by 1:
A A +
→ N+e +ν (30.3.6)
ZN Z−1
Electron capture has the same effect on the number of protons and neutrons in a nucleus as positron emission.
Alpha decay is the only type of radioactive decay that results in an appreciable change in an atom ‘s atomic mass. However, rather
than being destroyed, the two protons and two neutrons an atom loses in alpha decay are released as a helium nucleus.
Nuclear Fission
Chain reactions of nuclear fission release a tremendous amount of energy, but follow the Law of Conservation of Nucleon Number.
Consider, for example, the multistep reaction that occurs when a U-235 nucleus accepts a neutron, as in:
30.3.2 https://phys.libretexts.org/@go/page/16244
1 235 236 144 89 1
n+ U → U → Ba + Kr + 3 n + γ (30.3.7)
0 92 92 56 36 0
In each step, the total atomic mass of all species is a constant value of 236. This is the same with all fission reactions.
Nuclear Fusion
Finally, nuclear fusion follows the Law of Conservation of Nucleon Number. Consider the fusion of deuterium and tritium (both
hydrogen isotopes):
2 3 4 1
H+ H → He + n (30.3.8)
1 1 2 0
It is well understood that the tremendous amounts of energy released by nuclear fission and fusion can be attributed to the
conversion of mass to energy. However, the mass that is converted to energy is rather small compared to any sample, and never
includes the conversion of a proton or neutron to energy. Thus, the number of nucleons before and after fission and fusion is always
constant.
Key Points
Quantum tunneling applies to all objects facing any barrier. However, the probability of its occurrence is essentially negligible
for macroscopic purposes; it is only ever observed to any appreciable degree on the nanoscale level.
Quantum tunneling is explained by the imaginary component of the Schrödinger equation. Because the wave function of any
object contains an imaginary component, it can exist in imaginary space.
Tunneling decreases with the increasing mass of the object that must tunnel and with the increasing gap between the object’s
energy and the energy of the barrier it must overcome.
The law of Conservation of Nuclear Number states that the sum of protons and neutrons among species before and after a
nuclear reaction will be the same.
In radioactive decay, a proton can be converted to a neutron and a neutron can be converted to a proton (beta-decay).
Nuclear fusion and fission involve the conversion of matter to energy, but the matter that is converted is never a full nucleon.
Key Terms
tunneling: the quantum-mechanical passing of a particle through an energy barrier
fusion: A nuclear reaction in which nuclei combine to form more massive nuclei with the concomitant release of energy.
fission: The process of splitting the nucleus of an atom into smaller particles; nuclear fission.
nucleon: One of the subatomic particles of the atomic nucleus (i.e., a proton or a neutron).
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30.4: Applications of Nuclear Physics
Example 30.4.1:
The nuclear medicine whole body bone scan is generally used in evaluations of various bone related pathology, such as for
bone pain, stress fracture, nonmalignant bone lesions, bone infections, or the spread of cancer to the bone.
Radiation therapy involves the application of ionizing radiation to treat conditions such as hyperthyroidism, thyroid cancer, and
blood disorders. Radiation therapy is particularly effective as a treatment of a number of types of cancer if they are localized to one
area of the body. It may also be used as part of curative therapy, to prevent tumor recurrence after surgery, or to remove a primary
malignant tumor. Radiation therapy is synergistic with chemotherapy and has been used before, during, and after chemotherapy in
susceptible cancers.
Ionizing radiation works by damaging the DNA of exposed tissue, leading to cellular death. When external beam therapy is used,
shaped radiation beams are aimed from several angles of exposure to intersect at the tumor, providing a much larger absorbed dose
there than in the surrounding, healthy tissue.
External Beam Therapy: Radiation therapy of the pelvis. Lasers and a mold under the legs are used for precise positioning
Brachytherapy is another form of radiation therapy, in which a therapeutic radioisotope is injected into the body to chemically
localize to the tissue that requires destruction . A key feature of brachytherapy is that the irradiation affects only a very localized
area around the radiation sources. Exposure to radiation of healthy tissues further away from the sources is therefore reduced in this
technique.
Radiation therapy is in itself painless. Many low-dose palliative treatments (for example, radiation therapy targeting bony
metastases) cause minimal or no side effects, although short-term pain flare-ups can be experienced in the days following treatment
due to edemas compressing nerves in the treated area. Higher doses can cause varying side effects during treatment (acute), in the
months or years following treatment (long-term), or after re-treatment (cumulative). The nature, severity, and longevity of side
effects depend on the organs that receive the radiation, the treatment itself (type of radiation, dose, fractionation, concurrent
chemotherapy), and the individual patient.
Dosimetry
Radiation dosimetry is the measurement and calculation of the absorbed dose resulting from the exposure to ionizing radiation.
30.4.1 https://phys.libretexts.org/@go/page/16245
learning objectives
Explain difference between absorbed dose and dose equivalent.
Radiation dosimetry is the measurement and calculation of the absorbed dose in matter and tissue resulting from the exposure to
indirect and direct ionizing radiation.
Measuring Radiation
There are several ways of measuring the dose of ionizing radiation. Workers who come in contact with radioactive substances or
who may be exposed to radiation routinely carry personal dosimeters . In the United States, these dosimeters usually contain
materials that can be used in thermoluminescent dosimetry or optically stimulated luminescence. Outside the United States, the
most widely used type of personal dosimeter is the film badge dosimeter, which uses photographic emulsions that are sensitive to
ionizing radiation. The equipment used in radiotherapy (a linear particle accelerator in external beam therapy) is routinely
calibrated using ionization chambers or the new and more accurate diode technology. Internal dosimetry is used to evaluate the
intake of particles inside a human being.
Ionization Chamber: This ionization chamber was used in the South Atlantic Anomaly Probe project.
Dose is reported in grays (Gy) for absorbed doses or sieverts (Sv) for dose equivalents, where 1 Gy or 1 Sv is equal to 1 joule per
kilogram. Non-SI units are still prevalent as well: absorbed dose is often reported in rads and dose equivalent in rems. By
definition, 1 Gy = 100 rad, and 1 Sv = 100 rem.
Biological Effects
The distinction between absorbed dose (Gy/rad) and dose equivalent (Sv/rem) is based upon the biological effects. The weighting
factor (wr) and tissue/organ weighting factor (WT) have been established. They compare the relative biological effects of various
types of radiation and the susceptibility of different organs.
The weighting factor for the whole body is 1, such that 1 Gy of radiation delivered to the whole body is equal to one sievert.
Therefore, the WT for all organs in the whole body must sum to 1.
By definition, X-rays and gamma rays have a wr of unity, such that 1 Gy = 1 Sv (for whole-body irradiation). Values of wr are as
high as 20 for alpha particles and neutrons. That is to say, for the same absorbed dose in Gy, alpha particles are 20 times as
biologically potent as x-rays or gamma rays.
Dose is a measure of deposited dose and therefore can never decrease: removal of a radioactive source can reduce only the rate of
increase of absorbed dose — never the total absorbed dose.
30.4.2 https://phys.libretexts.org/@go/page/16245
learning objectives
Describe effects of ionizing radiation on living organisms.
Example 30.4.2:
The Radium Girls were female factory workers who contracted radiation poisoning from painting watch dials with glow-in-
the-dark paint at the United States Radium factory in Orange, New Jersey around 1917. The women, who had been told the
paint was harmless, ingested deadly amounts of radium by licking their paintbrushes to give them a fine point; some also
painted their fingernails and teeth with the glowing substance.
Ionizing radiation is generally harmful, even potentially lethal, to living organisms. Although radiation was discovered in the late
19th century, the dangers of radioactivity and of radiation were not immediately recognized. The acute effects of radiation were
first observed in the use of x-rays when Wilhelm Röntgen intentionally subjected his fingers to x-rays in 1895. The genetic effects
of radiation, including the effects on cancer risk, were recognized much later. In 1927, Hermann Joseph Muller published research
showing genetic effects.
Some effects of ionizing radiation on human health are stochastic, meaning that their probability of occurrence increases with dose,
while the severity is independent of dose. Radiation-induced cancer, teratogenesis, cognitive decline, and heart disease are all
examples of stochastic effects. Other conditions, such as radiation burns, acute radiation syndrome, chronic radiation syndrome,
and radiation-induced thyroiditis are deterministic, meaning they reliably occur above a threshold dose and their severity increases
with dose. Deterministic effects are not necessarily more or less serious than stochastic effects; either can ultimately lead to damage
ranging from a temporary nuisance to death.
30.4.3 https://phys.libretexts.org/@go/page/16245
Therapeutic Uses of Radiation
Radiation therapy uses ionizing radiation to treat conditions such as hyperthyroidism, cancer, and blood disorders.
learning objectives
Explain difference between external beam radiotherapy and brachytherapy.
Radiation therapy involves the application of ionizing radiation to treat conditions such as hyperthyroidism, thyroid cancer, and
blood disorders. Radiation therapy is particularly effective as a treatment of a number of types of cancer if they are localized to one
area of the body. It may also be used as part of curative therapy, to prevent tumor recurrence after surgery, or to remove a primary
malignant tumor. Radiation therapy is synergistic with chemotherapy and has been used before, during, and after chemotherapy in
susceptible cancers.
Ionizing radiation works by damaging the DNA of exposed tissue, leading to cellular death. When external beam therapy is used,
shaped radiation beams are aimed from several angles of exposure to intersect at the tumor, providing a much larger absorbed dose
there than in the surrounding, healthy tissue .
External Beam Therapy: Radiation therapy of the pelvis. Lasers and a mold under the legs are used for precise positioning
Brachytherapy is another form of radiation therapy, in which a therapeutic radioisotope is injected into the body to chemically
localize to the tissue that requires destruction . A key feature of brachytherapy is that the irradiation affects only a very localized
area around the radiation sources. Exposure to radiation of healthy tissues further away from the sources is therefore reduced in this
technique.
30.4.4 https://phys.libretexts.org/@go/page/16245
Clinical Applications of Brachytherapy: Body sites in which brachytherapy can be used to treat cancer
Radiation therapy is in itself painless. Many low-dose palliative treatments (for example, radiation therapy targeting bony
metastases) cause minimal or no side effects, although short-term pain flare-ups can be experienced in the days following treatment
due to edemas compressing nerves in the treated area. Higher doses can cause varying side effects during treatment (acute), in the
months or years following treatment (long-term), or after re-treatment (cumulative). The nature, severity, and longevity of side
effects depend on the organs that receive the radiation, the treatment itself (type of radiation, dose, fractionation, concurrent
chemotherapy), and the individual patient.
learning objectives
Explain how food irradiation is performed, commenting on its purpose and safety
Food irradiation is a process of treating a food to a specific dosage of ionizing radiation for a predefined length of time. This
process slows or halts spoilage that is due to the growth of pathogens. Food irradiation is currently permitted by over 50 countries,
and the volume of food treated is estimated to exceed 500,000 metric tons annually worldwide. Irradiated food is sold in regular
stores, often in specially marked packages .
Radura Logo: The Radura logo, required by U.S. Food and Drug Administration regulations to show a food has been treated with
ionizing radiation
By irradiating food, depending on the dose, some or all of the microorganisms, bacteria, viruses, and insects present are killed. This
prolongs the shelf-life of the food in cases where pathogenic spoilage is the limiting factor. Some foods, e.g., herbs and spices, are
irradiated at sufficient doses (five kilograys or more) to reduce the microbial counts by several orders of magnitude. Such
ingredients do not carry spoilage or pathogen microorganisms into the final product. It has also been shown that irradiation can
delay the ripening of fruits and the sprouting of vegetables.
Food irradiation using cobalt-60 is the preferred method by most processors. This is because the deep penetration of gamma rays
allows for the treatment of entire industrial pallets or totes at once, which reduces the need for material handling. A pallet or tote is
typically exposed for several minutes to several hours, depending on the dose. Radioactive material must be monitored and
carefully stored to shield workers and the environment from its gamma rays. During operation this is achieved using concrete
shields. With most designs the radioisotope can be lowered into a water-filled source storage pool to allow maintenance personnel
to enter the radiation shield. In this mode the water in the pool absorbs the radiation.
X-ray irradiators are considered an alternative to isotope-based irradiation systems. X-rays are generated by colliding accelerated
electrons with a dense material (the target), such as tantalum or tungsten, in a process known as bremsstrahlung-conversion. X-ray
irradiators are scalable and have deep penetration comparable to Co-60, with the added benefit that the electronic source stops
radiating when switched off. They also permit dose uniformity, but these systems generally have low energetic efficiency during
the conversion of electron energy to photon radiation, so they require much more electrical energy than other systems. X-ray
systems also rely on concrete shields to protect the environment and workers from radiation.
Irradiated food does not become radioactive, since the particles that transmit radiation are not themselves radioactive. Still, there is
some controversy in the application of irradiation due to its novelty, the association with the nuclear industry, and the potential for
the chemical changes to be different than the chemical changes due to heating food (since ionizing radiation produces a higher
energy transfer per collision than conventional radiant heat).
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Tracers
A radioactive tracer is a chemical compound in which one or more atoms have been replaced by a radioisotope.
learning objectives
Explain structure and use of radioactive tracers
A radioactive tracer is a chemical compound in which one or more atoms have been replaced by a radioisotope. By virtue of its
consequent radioactive decay, this compound can be used to explore the mechanism of chemical reactions by tracing the path that
the radioisotope follows from reactants to products.
The underlying principle in the creation of a radioactive tracer is that an atom in a chemical compound is replaced by another atom
of the same chemical element. In a tracer, this substituting atom is a radioactive isotope. This process is often called radioactive
labeling. Radioactive decay is much more energetic than chemical reactions. Therefore, the radioactive isotope can be present in
low concentration and its presence still detected by sensitive radiation detectors such as Geiger counters and scintillation counters.
Nuclear Fusion
In nuclear fusion two or more atomic nuclei collide at very high speed and join, forming a new nucleus.
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learning objectives
Analyze possibility of the use of nuclear fusion for the production of electricity.
Example 30.4.3:
The sun is a main-sequence star and therefore generates its energy through nuclear fusion of hydrogen nuclei into helium. In its
core, the sun fuses 620 million metric tons of hydrogen each second.
Nuclear fusion is a nuclear reaction in which two or more atomic nuclei collide at very high speed and join to form a new type of
atomic nucleus. During this process, matter is not conserved because some of the mass of the fusing nuclei is converted into energy.
Fusion of Deuterium with Tritium: Fusion of deuterium with tritium creating helium-4, freeing a neutron, and releasing 17.59
MeV of energy; some mass changes form to appear as the kinetic energy of the products
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Fission and Fusion
Describes the difference between fission and fusion
Fusion reactions of light elements power the stars and produce virtually all elements in a process called nucleosynthesis. The fusion
of lighter elements in stars releases energy and mass. For example, in the fusion of two hydrogen nuclei to form helium, 0.7 percent
of the mass is carried away from the system in the form of kinetic energy or other forms of energy (such as electromagnetic
radiation).
It takes considerable energy to force nuclei to fuse, even nuclei of the lightest element, hydrogen. This is because all nuclei have a
positive charge due to their protons, and since like charges repel, nuclei strongly resist being put close together. Accelerated to high
speeds, they can overcome this electrostatic repulsion and be forced close enough for the attractive nuclear force to be sufficiently
strong to achieve fusion. The fusion of lighter nuclei, which creates a heavier nucleus and often a free neutron or proton, generally
releases more energy than it takes to force the nuclei together. This is an exothermic process that can produce self-sustaining
reactions.
Research into controlled fusion, with the aim of producing fusion power for the production of electricity, has been conducted for
over 60 years. It has been accompanied by extreme scientific and technological difficulties, but it has resulted in progress. At
present, controlled fusion reactions have been unable to produce self-sustaining controlled fusion reactions. Researchers are
working on a reactor that theoretically will deliver 10 times more fusion energy than the amount needed to heat up plasma to
required temperatures. Workable designs of this reactor were originally scheduled to be operational in 2018; however, this has been
delayed, and a new date has not been released.
learning objectives
Explain how nuclear chain reactions can be controlled.
Example 30.4.4:
Some serious nuclear and radiation accidents have occurred. In 2011, three of the reactors at Fukushima I overheated, causing
meltdowns that eventually led to explosions, which released large amounts of radioactive material into the air.
Nuclear fission is a nuclear reaction in which the nucleus of an atom splits into smaller (lighter) nuclei. This reaction often
produces free neutrons and photons (in the form of gamma rays) and releases a very large amount of energy, even by the standards
of radioactive decay . The two nuclei produced are most often of comparable but slightly different sizes, typically with a mass ratio
of products of about 3 to 2, for common fissile nuclides.
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Nuclear Fission Reaction: An induced nuclear fission event. A neutron is absorbed by the nucleus of a uranium-235 atom, which
in turn splits into fast-moving lighter elements (fission products) and free neutrons
For example, when a large fissile atomic nucleus such as uranium-235 or plutonium-239 absorbs a neutron, it may undergo nuclear
fission. The heavy nucleus splits into two or more lighter nuclei (the fission products), releasing kinetic energy, gamma radiation,
and free neutrons. A portion of these neutrons may later be absorbed by other fissile atoms and trigger further fission events, which
release more neutrons, and so on. This is known as a nuclear chain reaction.
Just as conventional power stations generate electricity by harnessing the thermal energy released from burning fossil fuels, the
thermal energy released from nuclear fission can be converted in electricity by nuclear reactors. A nuclear chain reaction can be
controlled by using neutron poisons and neutron moderators to change the percentage of neutrons that will go on to cause more
fissions. Nuclear reactors generally have automatic and manual systems to shut the fission reaction down if unsafe conditions are
detected.
The reactor core generates heat in a number of ways. The kinetic energy of fission products is converted to thermal energy when
these nuclei collide with nearby atoms. Some of the gamma rays produced during fission are absorbed by the reactor, and their
energy is converted to heat. Heat is produced by the radioactive decay of fission products and materials that have been activated by
neutron absorption. This decay heat source will remain for some time even after the reactor is shut down.
A nuclear reactor coolant — usually water, but sometimes a gas, liquid metal, or molten salt — is circulated past the reactor core to
absorb the heat that it generates. The heat is carried away from the reactor and is then used to generate steam.
The power output of the reactor is adjusted by controlling how many neutrons are able to create more fissions. Control rods that are
made of a neutron poison are used to absorb neutrons. Absorbing more neutrons in a control rod means that there are fewer
neutrons available to cause fission, so pushing the control rod deeper into the reactor will reduce the reactor’s power output, and
extracting the control rod will increase it.
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Fukushima Daiichi Nuclear Disaster: Satellite image taken March 16, 2011 of the four damaged reactor buildings
Emission Topography
Positron emission tomography is a nuclear medical imaging technique that produces a three-dimensional image of processes in the
body.
learning objectives
Discuss possibility of uses of positron emission tomography with other diagnostic techniques.
Positron emission tomography (PET) is a nuclear medical imaging technique that produces a three-dimensional image or picture of
functional processes in the body. The system detects pairs of gamma rays emitted indirectly by a positron-emitting radionuclide
(tracer), which is introduced into the body on a biologically active molecule. Three-dimensional images of tracer concentration
within the body are then constructed by computer analysis.
PET acquisition process occurs as the radioisotope undergoes positron emission decay (also known as positive beta decay), it emits
a positron, an antiparticle of the electron with opposite charge. The emitted positron travels in tissue for a short distance (typically
less than 1 mm, but dependent on the isotope), during which time it loses kinetic energy, until it decelerates to a point where it can
interact with an electron. The encounter annihilates both electron and positron, producing a pair of annihilation (gamma) photons
moving in approximately opposite directions. These are detected when they reach a scintillator in the scanning device, creating a
burst of light which is detected by photomultiplier tubes or silicon avalanche photodiodes . The technique depends on simultaneous
or coincident detection of the pair of photons moving in approximately opposite directions (it would be exactly opposite in their
center of mass frame, but the scanner has no way to know this, and so has a built-in slight direction-error tolerance). Photons that
do not arrive in temporal “pairs” (i.e. within a timing-window of a few nanoseconds) are ignored.
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PET/CT-System: PET/CT-System with 16-slice CT; the ceiling mounted device is an injection pump for CT contrast agent.
PET scanning is non-invasive, but it does involve exposure to ionizing radiation. The total dose of radiation is significant, usually
around 5–7 mSv. However, in modern practice, a combined PET/CT scan is almost always performed, and for PET/CT scanning,
the radiation exposure may be substantial—around 23–26 mSv (for a 70 kg person—dose is likely to be higher for higher body
weights). When compared to the classification level for radiation workers in the UK of 6 mSv, it can be seen that use of a PET scan
needs proper justification.
Nuclear Weapons
A nuclear weapon is an explosive device that derives its destructive force from nuclear reactions—either fission, fusion, or a
combination.
learning objectives
Explain the difference between an “atomic” bomb and a “hydrogen” bomb, discussing their history
A nuclear weapon is an explosive device that derives its destructive force from nuclear reactions, either fission or a combination of
fission and fusion. Both reactions release vast quantities of energy from relatively small amounts of matter. The first fission (i.e.,
“atomic”) bomb test released the same amount of energy as approximately 20,000 tons of trinitrotoluene (TNT). The first fusion
(i.e., thermonuclear “hydrogen”) bomb test released the same amount of energy as approximately 10,000,000 tons of TNT.
A modern thermonuclear weapon weighing little more than 2,400 pounds (1,100 kg) can produce an explosive force comparable to
the detonation of more than 1.2 million tons (1.1 million tonnes) of TNT. Thus, even a small nuclear device no larger than
traditional bombs can devastate an entire city by blast, fire and radiation. Nuclear weapons are considered weapons of mass
destruction, and their use and control have been a major focus of international relations policy since their inception.
Only two nuclear weapons have been used in the course of warfare, both by the United States near the end of World War II. On
August 6, 1945, a uranium gun-type fission bomb code-named “Little Boy” was detonated over the Japanese city of Hiroshima.
Only three days later a plutonium implosion-type fission bomb code-named “Fat Man” (as illustrated in ) was exploded over
Nagasaki, Japan. The resulting mushroom cloud is shown in . The death toll from the two bombings was estimated at
approximately 200,000 people—mostly civilians, and mainly from acute injuries sustained from the explosions. The role of the
bombings in Japan’s surrender, and their ethical implications, remain the subject of scholarly and popular debate.
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Nagasaki Atomic Bombing: The mushroom cloud of the atomic bombing of Nagasaki, Japan (August 9,1945) rose some 18
kilometers (11 mi) above the bomb’s hypocenter.
Fat Man Atomic Bomb: The first nuclear weapons were gravity bombs, such as this “Fat Man” weapon dropped on Nagasaki,
Japan. They were very large and could only be delivered by heavy bomber aircraft.
Since the bombings of Hiroshima and Nagasaki, nuclear weapons have been detonated on over two thousand occasions for testing
purposes and demonstrations. Only a small number of nations either possess such weapons, or are suspected of trying to acquire
and/or develop them. The only countries known to have detonated nuclear weapons (and that acknowledge possessing such
weapons) are, as listed chronologically by date of first test: the United States, the Soviet Union (succeeded as a nuclear power by
Russia), the United Kingdom, France, the People’s Republic of China, India, Pakistan, and North Korea. In addition, it is also
widely believed that Israel possesses nuclear weapons (though they have not admitted to it).
The Federation of American Scientists estimates that as of 2012, there are more than 17,000 nuclear warheads in the world, with
around 4,300 considered “operational”—as in ready for use.
learning objectives
Explain the difference between magnetic resonance imaging and computed tomography.
Magnetic resonance imaging (MRI), also called nuclear magnetic resonance imaging (NMRI) or magnetic resonance tomography
(MRT), is a medical imaging technique used in radiology to visualize internal structures of the body in detail. MRI utilized the
property of nuclear magnetic resonance (NMR) to image the nuclei of atoms inside the body.
MRI machines (as pictured in ) make use of the fact that body tissue contains a large amount of water and therefore protons (1H
nuclei), which get aligned in a large magnetic field. Each water molecule has two hydrogen nuclei or protons. When a person is
inside the scanner’s powerful magnetic field, the hydrogen protons in their body align with the direction of the field. A radio
frequency current is briefly activated, producing a varying electromagnetic field. This electromagnetic field has just the right
frequency (known as the resonance frequency) to become absorbed and then reverse the rotation of the hydrogen protons in the
magnetic field.
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generated; this signal can be measured with receiver coils. Hydrogen protons in different tissues return to their equilibrium state at
different relaxation rates. Images are then constructed by performing a complex mathematical analysis of the signals emitted by the
hydrogen protons.
MRI shows a marked contrast between the different soft tissues of the body, making it especially useful in imaging the brain, the
muscles, the heart, and cancerous tissue—as compared with other medical imaging techniques such as computed tomography (CT)
or X-rays. MRI contrast agents may be injected intravenously to enhance the appearance of blood vessels, tumors or inflammation.
Unlike CT, MRI does not use ionizing radiation and is generally a very safe procedure. The strong magnetic fields and radio pulses
can, however, affect metal implants (including cochlear implants and cardiac pacemakers).
Key Points
Ionizing radiation works by damaging the DNA of exposed tissue, leading to cellular death.
In external beam radiotherapy, shaped radiation beams are aimed from several angles of exposure to intersect at the tumor,
providing a much larger absorbed dose there than in the surrounding, healthy tissue.
In brachytherapy, a therapeutic radioisotope is injected into the body to chemically localize to the tissue that requires
destruction.
There are several ways of measuring doses of ionizing radiation: personal dosimeters, ionization chambers, and internal
dosimetry.
The distinction between absorbed dose (Gy/rad) and dose equivalent (Sv/rem) is based upon the biological effects.
Dose is a measure of deposited dose and therefore can never decrease: removal of a radioactive source can reduce only the rate
of increase of absorbed dose — never the total absorbed dose.
The effects of ionizing radiation on human health are separated into stochastic effects (the probability of occurrence increases
with dose) and deterministic effects (they reliably occur above a threshold dose, and their severity increases with dose).
Quantitative data on the effects of ionizing radiation on human health are relatively limited compared to other medical
conditions because of the low number of cases to date and because of the stochastic nature of some of the effects.
Two pathways (external and internal) of exposure to ionizing radiation exist.
Ionizing radiation works by damaging the DNA of exposed tissue, leading to cellular death.
In external beam radiotherapy, shaped radiation beams are aimed from several angles of exposure to intersect at the tumor,
providing a much larger absorbed dose there than in the surrounding, healthy tissue.
In brachytherapy, a therapeutic radioisotope is injected into the body to chemically localize to the tissue that requires
destruction.
Food irradiation kills some of the microorganisms, bacteria, viruses, and insects found in food. It prolongs shelf-life in cases
where pathogenic spoilage is the limiting factor.
Food irradiation using cobalt-60 is the preferred method by most processors.
Irradiated food does not become radioactive, since the particles that transmit radiation are not themselves radioactive.
Radioactive tracers are used to explore the mechanism of chemical reactions by tracing the path that the radioisotope follows
from reactants to products.
The radioactive isotope can be present in low concentration and its presence still detected by sensitive radiation detectors.
All the commonly used radioisotopes have short half-lives, do not occur in nature, and are produced through nuclear reactions.
The fusion of lighter elements releases energy.
Matter is not conserved during fusion reactions.
Fusion reactions power the stars and produce virtually all elements in a process called nucleosynthesis.
Nuclear fission is a nuclear reaction in which the nucleus of an atom splits into smaller parts, releasing a very large amount of
energy.
Nuclear chain reactions can be controlled using neutron poisons and neutron moderators.
Although the nuclear power industry has improved the safety and performance of reactors and has proposed new, safer reactor
designs, there is no guarantee that serious nuclear accidents will not occur.
PET scanning utilizes detection of gamma rays emitted indirectly by a positron-emitting radionuclide (tracer), which is
introduced into the body on a biologically active molecule.
PET scans are increasingly read alongside CT or magnetic resonance imaging (MRI) scans, with the combination giving both
anatomic and metabolic information.
PET scanning is non-invasive, but it does involve exposure to ionizing radiation.
Nuclear weapons utilize either fission (“atomic” bomb) or combination of fission and fusion (“hydrogen” bomb).
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Nuclear weapons are considered weapons of mass destruction.
The use and control of nuclear weapons is a major focus of international relations policy since their first use.
MRI makes use of the property of nuclear magnetic resonance to image nuclei of atoms inside the body.
MRI provides good contrast between the different soft tissues of the body (making it especially useful in imaging the brain, the
muscles, the heart, and cancerous tissue).
Although MRI uses non-ionizing radiation, the strong magnetic fields and radio pulses can affect metal implants, including
cochlear implants and cardiac pacemakers.
Key Terms
external beam therapy: Radiotherapy that directs the radiation at the tumour from outside the body.
ionizing radiation: high-energy radiation that is capable of causing ionization in substances through which it passes; also
includes high-energy particles
brachytherapy: Radiotherapy using radioactive sources positioned within (or close to) the treatment volume.
diode: an electronic device that allows current to flow in one direction only; a valve
dosimeter: A dosimeter is a device used to measure a dose of ionizing radiation. These normally take the form of either
optically stimulated luminescence (OSL), photographic-film, thermoluminescent (TLD), or electronic personal dosimeters
(PDM).
gamma ray: A very high frequency (and therefore very high energy) electromagnetic radiation emitted as a consequence of
radioactivity.
x-ray: Short-wavelength electromagnetic radiation usually produced by bombarding a metal target in a vacuum. Used to create
images of the internal structure of objects; this is possible because x-rays pass through most objects and can expose
photographic film
radioactive tracer: a radioactive isotope that, when injected into a chemically similar substance, or artificially attached to a
biological or physical system, can be traced by radiation detection devices
isotope: any of two or more forms of an element where the atoms have the same number of protons but a different number of
neutrons within their nuclei. As a consequence, atoms for the same isotope will have the same atomic number but different mass
numbers (atomic weights)
radioactive decay: any of several processes by which unstable nuclei emit subatomic particles and/or ionizing radiation and
disintegrate into one or more smaller nuclei
nucleosynthesis: any of several processes that lead to the synthesis of heavier atomic nuclei
fusion: A nuclear reaction in which nuclei combine to form more massive nuclei with the concomitant release of energy.
control rod: any of a number of steel tubes, containing boron or another neutron absorber, that is inserted into the core of a
nuclear reactor in order to control its rate of reaction
nuclear reactor: any device in which a controlled chain reaction is maintained for the purpose of creating heat (for power
generation) or for creating neutrons and other fission products for experimental, medical, or other purposes
fission: The process of splitting the nucleus of an atom into smaller particles; nuclear fission.
tracer: A chemical used to track the progress or history of a natural process.
positron: The antimatter equivalent of an electron, having the same mass but a positive charge.
tomography: Imaging by sections or sectioning.
warfare: The waging of war or armed conflict against an enemy.
computed tomography: (CT) – A form of radiography which uses computer software to create images, or slices, at various
planes of depth from images taken around a body or volume of interest.
nuclear magnetic resonance: (NMRI) – The absorption of electromagnetic radiation (radio waves), at a specific frequency, by
an atomic nucleus placed in a strong magnetic field; used in spectroscopy and in magnetic resonance imaging.
magnetic resonance imaging: Commonly referred to as MRI; a technique that uses nuclear magnetic resonance to form cross
sectional images of the human body for diagnostic purposes.
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CHAPTER OVERVIEW
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1
4.1: Introduction
Newton and His Laws
There are three laws of motion that describe the relationship between forces, mass, and acceleration.
learning objectives
Apply three Newton’s laws of motion to relate forces, mass, and acceleration
Newton’s laws of motion describe the relationship between the forces acting on a body and its motion due to those forces. For
example, if your car breaks down and you need to push it, you must exert a force with your hands on the car in order for it to move.
The laws of motion will tell you how quickly the car will move from your pushing. There are three laws of motion:
First law: If an object experiences no net force, then its velocity is constant: the object is either at rest (if its velocity is zero), or it
moves in a straight line with constant speed (if its velocity is nonzero). For example, if you don’t push the car (no force), then it
doesn’t move.
Second law: The acceleration aa of a body is parallel and directly proportional to the net force F acting on the body, is in the
direction of the net force, and is inversely proportional to the mass mm of the body:
F
F = m ⋅ a or a = (4.1.1)
m
For example, if you push the car with a greater force it will accelerate more. But, if the car is more massive (mm is larger) then it
won’t accelerate as much from the same size force as a lighter car.
Third law: When a first body exerts a force F on a second body, the second body simultaneously exerts a force F = −F on the
1 2 1
first body. This means that F and F are equal in magnitude and opposite in direction. For example, when you push a car, if it is
1 2
exerting the same force on you that you are exerting on it, you might wonder why you don’t move backwards? The answer is there
are also forces from the ground on your feet pushing you forward. So, in fact, the car is pushing a force back on you that is of the
same magnitude that you are using to push it forward.
In the figure below there are some practical examples illustrating the concept of force:
Strain: by using a machine known as pulley you can easily raise or lower a massive body
Gravitational Force: a massive body is attracted downward by the gravitational force practiced by the Earth
Magnetic Force: two magnets repel each other when the same poles get closer
Key Points
Acceleration of an object is proportional to the force on it.
Force causes an object to move.
Objects with more mass require more force to move.
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Key Terms
force: Any influence that causes an object to undergo a certain change, either concerning its movement, direction or
geometrical construction.
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4.2: Force and Mass
Force
Force is any influence that causes an object to change, either concerning its movement, direction, or geometrical construction.
learning objectives
Develop the relationship between mass and acceleration in determining force
Overview of Forces
In physics, a force is any influence that causes an object to undergo a certain change, either concerning its movement, direction, or
geometrical construction. It is measured with the SI unit of Newtons. A force is that which can cause an object with mass to change
its velocity, i.e., to accelerate, or which can cause a flexible object to deform. Force can also be described by intuitive concepts
such as a push or pull. A force has both magnitude and direction, making it a vector quantity.
What is a force?
What is a force?: Describes what forces are and what they do.
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Qualities of Force
The original form of Newton’s second law states that the net force acting upon an object is equal to the rate at which its momentum
changes. This law is further given to mean that the acceleration of an object is directly proportional to the net force acting on the
object, is in the direction of the net force, and is inversely proportional to the mass of the object.
As we mentioned, force is a vector quantity. A vector is a one dimensional array with elements of both magnitude and direction. In
a force vector, the mass, m, is the magnitude component and the acceleration, a, is the directional component. The equation for
force is written:
F = m⋅a (4.2.1)
Related concepts to force include thrust, which increases the velocity of an object; drag, which decreases the velocity of an object;
and torque which produces changes in rotational speed of an object. Forces which do not act uniformly on all parts of a body will
also cause mechanical stresses, a technical term for influences which cause deformation of matter. While mechanical stress can
remain embedded in a solid object, gradually deforming it, mechanical stress in a fluid determines changes in its pressure and
volume.
Dynamics
Dynamics is the study of the forces that cause objects and systems to move. To understand this, we need a working definition of
force. Our intuitive definition of force — that is, a push or a pull —is a good place to start. We know that a push or pull has both
magnitude and direction (therefore, it is a vector quantity) and can vary considerably in each regard.
shows a few examples of the “push-pull” nature of force. The top left example is that of a pulley system. The force that someone
would have to pull down on the cable would have to equal and exceed the force made by the mass the object and the effects of
gravity on those object in order for the system to move up. The top right example shows that any object resting on a surface will
still exert force on that surface. The bottom example is that of two magnets being attracted to each other due to magnetic force.
Mass
Mass is a physical property of matter that depends on size and shape of matter, and is expressed as kilograms by the SI system.
learning objectives
Justify the significance of understanding mass in physics
What is Mass?
All elements have physical properties whose values can help describe an elements physical state. Changes to these properties can
describe elemental transformations. Physical properties do not change the chemical nature of matter. The physical property we are
covering in this atom is called mass.
Mass is defined as a quantitative measure of an object’s resistance to acceleration. The terms mass and weight are often
interchanged, however it is incorrect to do so. Weight is a different property of matter that, while related to mass, is not mass, but
rather the amount of gravitational force acting on a given body of matter. Mass is an intrinsic property that never changes.
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Units of Mass
In order to measure something, a standard value must be established to use in relation to the object of measurement. This relation is
called a unit. The International System of Units (SI) measures mass in kilograms, or kg. There are other units of mass, including the
following (only the first two are accepted by the SI system):
t – Tonne; 1t = 1000kg
u – atomic mass unit; 1u ~= 1.66×10-27kg
sl – slug
lb – pound
2
2
Key Points
Force is stated as a vector quantity, meaning it has elements of both magnitude and direction. Mass and acceleration
respectively.
In layman’s terms, force is a push or pull that can be defined in terms of various standards.
Dynamics is the study of the force that causes objects and systems to move or deform.
External forces are any outside forces that act on a body, and internal forces are any force acting within a body.
Mass is defined as a quantitative measure of an object’s resistance to acceleration.
According to Newton’s second law of motion, if a body of fixed mass m is subjected to a single force F, its acceleration a is
given by F/m.
Mass is central in many concepts of physics, including:weight, momentum, acceleration, and kinetic energy.
According to Newton’s second law of motion, if a body of fixed mass m is subjected to a single force F, its acceleration a is
given by F/m.
Key Terms
force: A force is any influence that causes an object to undergo a certain change, either concerning its movement, direction or
geometrical construction.
velocity: A vector quantity that denotes the rate of change of position with respect to time, or a speed with a directional
component.
vector: A directed quantity, one with both magnitude and direction; the between two points.
mass: The quantity of matter which a body contains, irrespective of its bulk or volume. It is one of four fundamental properties
of matter. It is measured in kilograms in the SI system of measurement.
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4.3: Newton’s Laws
The First Law: Inertia
Newton’s first law of motion describes inertia. According to this law, a body at rest tends to stay at rest, and a body in motion tends
to stay in motion, unless acted on by a net external force.
learning objectives
Define the First Law of Motion
History
Sir Isaac Newton was an English scientist who was interested in the motion of objects under various conditions. In 1687, he
published a work called Philosophiae Naturalis Principla Mathematica, which described his three laws of motion. Newton used
these laws to explain and explore the motion of physical objects and systems. These laws form the basis for mechanics. The laws
describe the relationship between forces acting on a body and the motions experienced due to these forces. The three laws are as
follows:
1. If an object experiences no net force, its velocity will remain constant. The object is either at rest and the velocity is zero or it
moves in a straight line with a constant speed.
2. The acceleration of an object is parallel and directly proportional to the net force acting on the object, is in the direction of the
net force, and is inversely proportional to the mass of the object.
3. When a first object exerts a force on a second object, the second object simultaneously exerts a force on the first object,
meaning that the force of the first object and the force of the second object are equal in magnitude and opposite in direction.
Example 4.3.1:
If you are ice skating, and you push yourself away from the side of the rink, according to Newton’s first law you will continue
all the way to the other side of the rink. But, this won’t actually happen. Newton says that a body in motion will stay in motion
until an outside force acts upon it. In this and most other real world cases, this outside force is friction. The friction between
your ice skates and the ice is what causes you to slow down and eventually stop.
Let’s look at another situation. Refer to for this example. Why do we wear seat belts? Obviously, they’re there to protect us
from injury in case of a car accident. If a car is traveling at 60 mph, the driver is also traveling at 60 mph. When the car
suddenly stops, an external force is applied to the car that causes it to slow down. But there is no force acting on the driver, so
the driver continues to travel at 60 mph. The seat belt is there to counteract this and act as that external force to slow the driver
down along with the car, preventing them from being harmed.
Newton’s First Law: Newton’s first law in effect on the driver of a car
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Inertia
Sometimes this first law of motion is referred to as the law of inertia. Inertia is the property of a body to remain at rest or to remain
in motion with constant velocity. Some objects have more inertia than others because the inertia of an object is equivalent to its
mass. This is why it is more difficult to change the direction of a boulder than a baseball.
Doc Physics – Newton: Newton’s first law is hugely counterintuitive. You may have learned it in gradeschool, though. Let’s see it
for the mind-blowing conclusion it really is.
learning objectives
Define the Second Law of Motion
English scientist Sir Isaac Newton examined the motion of physical objects and systems under various conditions. In 1687, he
published his three laws of motion in Philosophiae Naturalis Principla Mathematica. The laws form the basis for mechanics—they
describe the relationship between forces acting on a body, and the motion experienced due to these forces. These three laws state:
1. If an object experiences no net force, its velocity will remain constant. The object is either at rest and the velocity is zero, or it
moves in a straight line with a constant speed.
2. The acceleration of an object is parallel and directly proportional to the net force acting on the object, is in the direction of the
net force and is inversely proportional to the mass of the object.
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3. When a first object exerts a force on a second object, the second object simultaneously exerts a force on the first object,
meaning that the force of the first object and the force of the second object are equal in magnitude and opposite in direction.
The first law of motion defines only the natural state of the motion of the body (i.e., when the net force is zero). It does not allow us
to quantify the force and acceleration of a body. The acceleration is the rate of change in velocity; it is caused only by an external
force acting on it. The second law of motion states that the net force on an object is equal to the rate of change of its linear
momentum.
Linear Momentum
Linear momentum of an object is a vector quantity that has both magnitude and direction. It is the product of mass and velocity of a
particle at a given time:
p = mv (4.3.1)
where, p = momentum, m = mass, and v = velocity . From this equation, we see that objects with more mass will have more
momentum.
Force and Mass: This animation demonstrates the connection between force and mass.
dp
F = (4.3.2)
dt
d(m ⋅ v)
F = (4.3.3)
dt
where, F = Force and t = time . From this we can further simplify the equation:
d(v)
F =m (4.3.4)
dt
F = m⋅a (4.3.5)
where, a=accelerationa=acceleration. As we stated earlier, acceleration is the rate of change of velocity, or velocity divided by time.
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Newton's Second Law (1 of 2) | Newton's Three Law…
Law…
Newton’s Three Laws of Mechanics – Second Law – Part 1: Here we’ll see how many people can confuse your understanding of
Newton’s 2nd law of motion through oversight, sloppy language, or cruel intentions.
4.3.4 https://phys.libretexts.org/@go/page/14454
Newton’s Three Laws of Mechanics – Second Law – Part Two: Equilibrium is investigated and Newton’s 1st law is seen as a
special case of Newton’s 2nd law!
learning objectives
Define the Third Law of Motion
Sir Isaac Newton was a scientist from England who was interested in the motion of objects under various conditions. In 1687, he
published a work called Philosophiae Naturalis Principla Mathematica, which contained his three laws of motion. Newton used
these laws to explain and explore the motion of physical objects and systems. These laws form the bases for mechanics. The laws
describe the relationship between forces acting on a body, and the motion is an experience due to these forces. Newton’s three laws
are:
1. If an object experiences no net force, its velocity will remain constant. The object is either at rest and the velocity is zero or it
moves in a straight line with a constant speed.
2. The acceleration of an object is parallel and directly proportional to the net force acting on the object, is in the direction of the
net force and is inversely proportional to the mass of the object.
3. When a first object exerts a force on a second object, the second object simultaneously exerts a force on the first object,
meaning that the force of the first object and the force of the second object are equal in magnitude and opposite in direction.
FA = −FB (4.3.6)
In this example, FA is the action and FB is the reaction. You have undoubtedly witnessed this law of motion. For example, take a
swimmer who uses her feet to push off the wall in order to gain speed. The more force she exerts on the wall, the harder she pushes
off. This is because the wall exerts the same force on her that she forces on it. She pushes the wall in the direction behind her,
therefore the wall will exert a force on her that is in the direction in front of her and propel her forward.
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Newton’s Third Law of Motion: When a swimmer pushes off the wall, the swimmer is using the third law of motion.
Take as another example, the concept of thrust. When a rocket launches into outer space, it expels gas backward at a high velocity.
The rocket exerts a large backward force on the gas, and the gas exerts and equal and opposite reaction force forward on the rocket,
causing it to launch. This force is called thrust. Thrust is used in cars and planes as well.
Newton’s Third Law: The most fundamental statement of basic physical reality is also the most often misunderstood. As your
mom if she’s clear on Newton’s Third. Then ask her why things can move if every force has a paired opposite force all the time,
forever.
Key Points
Newton’s three laws of physics are the basis for mechanics.
The first law states that a body at rest will stay at rest until a net external force acts upon it and that a body in motion will
remain in motion at a constant velocity until acted on by a net external force.
Net external force is the sum of all of the forcing acting on an object.
Just because there are forces acting on an object doesn’t necessarily mean that there is a net external force; forces that are equal
in magnitude but acting in opposite directions can cancel one another out.
Friction is the force between an object in motion and the surface on which it moves. Friction is the external force that acts on
objects and causes them to slow down when no other external force acts upon them.
Inertia is the tendency of a body in motion to remain in motion. Inertia is dependent on mass, which is why it is harder to
change the direction of a heavy body in motion than it is to change the direction of a lighter object in motion.
Newton’s three laws of motion explain the relationship between forces acting on an object and the motion they experience due
to these forces. These laws act as the basis for mechanics.
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The second law explains the relationship between force and motion, as opposed to velocity and motion. It uses the concept of
linear momentum to do this.
Linear momentum p, is the product of mass m, and velocity v : p = mv .
The second law states that the net force is equal to the derivative, or rate of change of its linear momentum.
By simplifying this relationship and remembering that acceleration is the rate of change of velocity, we can see that the second
law of motion is where the relationship between force and acceleration comes from.
If an object A exerts a force on object B, object B exerts an equal and opposite force on object A.
Newton’s third law can be seen in many everyday circumstances. When you walk, the force you use to push off the ground
backwards makes you move forward.
Thrust is an application of the third law of motion. A helicopter uses thrust to push the air under the propeller down, and
therefore lift off the ground.
Key Terms
inertia: The property of a body that resists any change to its uniform motion; equivalent to its mass.
friction: A force that resists the relative motion or tendency to such motion of two bodies in contact.
uniform motion: Motion at a constant velocity (with zero acceleration). Note that an object in motion will not change its
velocity unless an unbalanced force acts upon it.
net force: The combination of all the forces that act on an object.
momentum: (of a body in motion) the product of its mass and velocity.
acceleration: The amount by which a speed or velocity increases (and so a scalar quantity or a vector quantity).
symmetry: Exact correspondence on either side of a dividing line, plane, center or axis.
thrust: The force generated by propulsion, as in a jet engine.
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4.4: Other Examples of Forces
Weight
Weight is taken as the force on an object due to gravity, and is different than the mass of an object.
learning objectives
Infer what factors other than gravity will contribute to the apparent weight of an object
In physics, it is important to differentiate the weight of an object from its mass. The mass of an object is an intrinsic quantity,
independent of the location of the object. On the other hand, the weight of an object is an extrinsic quantity. It is considered as the
force on an object due to gravity. Since gravitational acceleration changes depending on the location in the universe, weight does as
well.
Mathematically, the weight of an object (W) can be found by multiplying its mass (m) by the acceleration due to gravity (g):
W = M ⋅ g . The strength of gravity varies very little over the surface of the Earth. In fact, the greatest percent difference in the
s2
.
The weight of an object has the same SI unit as force—the Newton (1N = 1kg ⋅ m
2
s
).
In US customary units, the weight of an object can be expressed in pounds. Keep in mind that in US units the pound is either a unit
of force or of mass. If one must find the weight (as opposed to the mass) of an object in US units, it can be calculated in terms of
pounds of force.
It is important to note that the apparent weight of an object (i.e., the weight of an object determined by a scale) will vary if forces
other than gravity are acting upon the object. For example, if you weigh a given mass underwater you will find a different result
than if you weigh that mass in air. In this case, the weight of the object varies due to the force of buoyancy. While the mass is in the
water it displaces fluid, resulting in an upward force upon it. This upward force affects the net force that the mass exerts on the
scale, and thus alters its “apparent” weight.
Spring Scale: A spring scale measures weight by finding the extent to which a spring is compressed.This is proportional to the
force that a mass exerts on the scale due to its weight.
4.4.1 https://phys.libretexts.org/@go/page/14455
Normal Forces
The normal force comes about when an object contacts a surface; the resulting force is always perpendicular to the surface of
contact.
learning objectives
Evaluate Newton’s Second and Third Laws in determining the normal force on an object
Overview
The normal force, F , comes about when an object contacts a surface. According to Newton’s third law, when one object exerts a
N
force on a second object, the second object always exerts a force that is equal in magnitude and opposite in direction on the first
object. This is the reason that the normal force exists.
A common situation in which a normal force exists is when a person stands on the ground. Because of Newton’s third law, the
ground exerts a force on the person that is equal in magnitude to the person’s weight. In this simple case, the weight of the person
and the opposing normal force are the only two forces considered on the person. The person remains still because the forces due to
weight and the normal force create a net force of zero on the person.
where:
m is the mass under consideration,
g is the acceleration due to gravity,
and θ is the angle between the inclined surface and the horizontal.
Inclined Plane: A mass rests on an inclined plane that is at an angle θ to the horizontal. The following forces act on the mass: the
weight of the mass (m ⋅ g ),the force due to friction (Fr),and the normal force (Fn).
Another interesting example involving normal forces is when a person stands in an elevator. When the elevator goes up, the normal
force is actually greater than the force due to gravity. In this situation there are only two forces acting on the person. The first is the
force of gravity on the person, which does not change. The second is the normal force. By summing the forces and setting them
equal to m⋅am⋅a (utilizing Newton’s second law), we find:
FN − m ⋅ g = m ⋅ a (4.4.2)
where:
FN is the normal force,
m⋅g is the force due to gravity,
m is the mass of the person,
4.4.2 https://phys.libretexts.org/@go/page/14455
Key Points
Weight is taken to be the force on an object due to gravity.
Weight and mass are not the same thing!
The weight of a given mass will be different when the acceleration due to gravity is different.
Apparent weight can change because of the effect of buoyancy.
The strength of gravity is almost the same everywhere on the surface of the Earth.
The normal force, F , comes about when an object contacts a surface.
N
The normal force exists because for every force, there is always an equal and opposite force.
The normal force is always perpendicular to the plane that the object contacts or rests on.
Key Terms
Gravitational acceleration: Gravitational acceleration is the acceleration that an object undergoes due solely to gravity
perpendicular: at or forming a right angle (to).
normal: A line or vector that is perpendicular to another line, surface, or plane.
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4.5: Problem-Solving
A General Approach
Basic problem-solving techniques can aid in the solution of problems involving motion (i.e., the laws of motion).
learning objectives
Assess the laws of motion through practiced problem solving techniques
When dealing with the laws of motion, although knowledge of concepts and equations is important, understanding basic problem
solving techniques can simplify the process of solving problems that may appear difficult. Your approach to problem solving can
involve several key steps.
Free body diagram: An example of a drawing to help identify forces and directions.
First, gather all relevant information from the problem. Identify all quantities that are given (the knowns), then do the same for all
quantities needed (the unknowns). Also, identify the physical principles involved (e.g., force, gravity, friction, etc. ).
Next, a drawing may be helpful. Sometimes a drawing can even help determine the known and unknown quantities. It need not be a
work of art, but it should be clear enough to illustrate proper dimension, (meaning one, two, or three dimensions). You can then use
this drawing to determine which direction is positive and which is negative (making note of this on the drawing).
A next step is to use what is known to find the appropriate equation to find what is unknown. While it is easiest to find an equation
that leaves only one unknown, sometimes this is not possible. In these situations, you can solve multiple equations to find the right
answer. Remember that equations represent physical principles and relationships, so use the equations and drawings in tandem.
You may then substitute the knowns into the appropriate equations and find a numerical solution.
Check the answer to see if it is reasonable and makes sense. Your judgment will improve and fine tune as you solve more problems
of this nature. This “judgement” step helps intuit the problem in terms of its conceptual meaning. If you can judge whether the
answer is reasonable, you have a deeper understanding of physics than simply the mechanics of problem solving.
When solving problems, we tend to perform these steps in different order, as well as do several steps simultaneously. There is no
rigid procedure that will work every time. Creativity and insight grow with experience. In time, the basics of problem solving can
become relatively automatic.
Key Points
Gathering all relevant information and identifying knowns and unknowns is an important first step.
Always make a drawing to help identify directions of forces and to establish x, y, and z axes.
Choose the correct equations, solve the problem, and check that the answer fits expectations numerically.
4.5.1 https://phys.libretexts.org/@go/page/14456
Key Terms
equation: An assertion that two expressions are equal, expressed by writing the two expressions separated by an equal sign;
from which one is to determine a particular quantity.
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4.6: Vector Nature of Forces
Forces in Two Dimensions
Forces act in a particular direction and have sizes dependent upon how strong the push or pull is.
learning objectives
Explain why forces are classified as “vector quantities”
Forces act in a particular direction and have sizes dependent upon how strong the push or pull is. Because of these characteristics,
forces are classified as “vector quantities. ” This means that forces follow a different set of mathematical rules than physical
quantities that do not have direction (denoted scalar quantities).
For example, when determining what happens when two forces act on the same object, it is necessary to know both the magnitude
and the direction of both forces to calculate the result. If both of these pieces of information are not known for each force, the
situation is ambiguous. For example, if you know that two people are pulling on the same rope with known magnitudes of force but
you do not know which direction either person is pulling, it is impossible to determine what the acceleration of the rope will be.
The two people could be pulling against each other as in tug of war or the two people could be pulling in the same direction. In this
simple one-dimensional example, without knowing the direction of the forces it is impossible to decide whether the net force is the
result of adding the two force magnitudes or subtracting one from the other. Associating forces with vectors avoids such problems.
When two forces act on a point particle, the resulting force or the resultant (also called the net force) can be determined by
following the parallelogram rule of vector addition: the addition of two vectors represented by sides of a parallelogram gives an
equivalent resultant vector which is equal in magnitude and direction to the transversal of the parallelogram. The magnitude of the
resultant varies from the difference of the magnitudes of the two forces to their sum, depending on the angle between their lines of
action.
Free-body diagrams can be used as a convenient way to keep track of forces acting on a system. Ideally, these diagrams are drawn
with the angles and relative magnitudes of the force vectors preserved so that graphical vector addition can be done to determine
the net force.
4.6.1 https://phys.libretexts.org/@go/page/14457
Forces as Vectors: Free-body diagrams of an object on a flat surface and an inclined plane. Forces are resolved and added together
to determine their magnitudes and the net force.
Key Points
When determining what happens when two forces act on the same object, it is necessary to know both the magnitude and the
direction of both forces to calculate the result.
When two forces act on a point particle, the resulting force or the resultant (also called the net force ), can be determined by
following the parallelogram rule of vector addition.
Free-body diagrams can be used as a convenient way to keep track of forces acting on an object.
Key Terms
vector: A directed quantity, one with both magnitude and direction; the between two points.
free-body diagram: A free body diagram, also called a force diagram, is a pictorial representation often used by physicists and
engineers to analyze the forces acting on a body of interest.
resultant: A vector that is the vector sum of multiple vectors
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4.7: Further Applications of Newton’s Laws
Applications of Newton’s Laws
Net force affects the motion, postion and/or shape of objects (some important and commonly used forces are friction, drag and
deformation).
learning objectives
Explain the effect of forces on an object’s motion and shape
We know that a net force affects the motion, position and shape of an object. It is useful at this point to look at some particularly
interesting and common forces that will provide further applications of Newton’s laws of motion. Specifically, we will discuss the
forces of friction, air or liquid drag, and deformation.
Friction
Friction is a force that resists movement between two surfaces sliding against each other. When surfaces in contact move relative to
each other, the friction between the two surfaces converts kinetic energy into heat. This property can have a dramatic effect, as seen
in the use of friction created by rubbing pieces of wood together to start a fire. Friction is not itself a fundamental force, but arises
from fundamental electromagnetic forces between the charged particles constituting the two contacting surfaces.
Drag
Another interesting force in everyday life is the force of drag on an object when it is moving in a fluid (either gas or liquid). You
feel this drag force when you move your hand through water, or through the wind. Like friction, the force of drag is a force that
resists motion. As we will discuss in later units, the drag force is proportional to the velocity of the object moving through it. We
see an illustrated example of drag force in.
Drag Force on a Barge: (a) A view from above of two tugboats pushing on a barge. (b) The free-body diagram for the ship
contains only forces acting in the plane of the water. It omits the two vertical forces—the weight of the barge and the buoyant force
of the water supporting it cancel and are not shown. Since the applied forces are perpendicular, the x – and y -axes are in the same
direction as F and F . The problem quickly becomes a one-dimensional problem along the direction of F , since friction is in
x y app
Deformation
We now move from consideration of forces that affect the motion of an object (such as friction and drag) to those that affect an
object’s shape. If a bulldozer pushes a car into a wall, the car will not move but it will noticeably change shape. The change in
shape of an object due to the application of a force is a deformation. Even very small forces are known to cause some deformation.
For small deformations, two important characteristics are observed. First, the object returns to its original shape when the force is
removed (that is, the deformation is elastic for small deformations). Second, the size of the deformation is proportional to the force.
Friction: Kinetic
If two systems are in contact and moving relative to one another, then the friction between them is called kinetic friction.
4.7.1 https://phys.libretexts.org/@go/page/15229
learning objectives
Explain the dynamics of energy for friction between two surfaces
When surfaces in contact move relative to each other, the friction between the two surfaces converts kinetic energy into heat. This
property can have dramatic consequences, as illustrated by the use of friction created by rubbing pieces of wood together to start a
fire. Kinetic energy is converted to heat whenever motion with friction occurs, for example when a viscous fluid is stirred.
Kinetic Friction Introduction: Here, I’ll explain the microscopic justification of friction and what we can know about it. The
coefficient of friction, too!
Kinetic (or dynamic) friction occurs when two objects are moving relative to each other and rub together; a sled on the ground
would be a good example of kinetic friction.
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Friction: Frictional forces always oppose motion or attempted motion between objects in contact. Friction arises in part because of
the roughness of the surfaces in contact, as seen in the expanded view. In order for the object to move, it must rise to where the
peaks can skip along the bottom surface. Thus, a force is required just to set the object in motion. Some of the peaks will be broken
off, also requiring a force to maintain motion. Much of the friction is actually due to attractive forces between molecules making up
the two objects, so that even perfectly smooth surfaces are not friction-free. Such adhesive forces also depend on the substances the
surfaces are made of, explaining, for example, why rubber-soled shoes slip less than those with leather soles.
The force of friction is what slows an object sliding over a surface. This force is what makes the brakes on cars work or causes
resistance when you slide your hand across a surface. The force of friction can be represented by an equation: F = μF . In
f riction n
this equation μ is something called the coefficient of friction. This is a unitless number that represents the strength of the friction of
the object. A very “grippy” surface like rubber might have a high coefficient of friction, whereas a slippery surface like ice has a
much lower coefficient. F is called the normal force and is the force of the surface pushing up on the object. In most cases on
n
level ground, the normal force will be the equal and opposite of the object’s weight. In other words, it is the force that the surface
must exert to keep the object from falling through.
The coefficient of kinetic friction is typically represented as μ and is usually less than the coefficient of static friction for the same
k
materials.
Friction: Static
Static friction is a type of friction that occurs to resist motion when two objects are at rest against each other.
learning objectives
Demonstrate the relationship of maximum force of static friction
Static Friction
Another type of frictional force is static friction, otherwise known as stiction. Like all friction, it acts to resist the motion of an
object moving over a surface. Unlike kinetic friction, however, static friction acts to resist the start of motion.
4.7.3 https://phys.libretexts.org/@go/page/15229
Static Friction and some friction challenges: Here, I talk about sneaky ol’ static friction.
Static friction is friction between two objects that are not moving relative to each other. This frictional force is what prevents a
parked car from sliding down a hill, for example. Before an object at rest on a surface can move, it must overcome the force of
static friction.
Static friction originates from multiple sources. For any given material on another material of the same composition, friction will
be greater as the material surfaces become rougher (consider sandpaper) on the macroscopic level. Additionally, intermolecular
forces can greatly influence friction when two materials are put into contact. When surface area is below the micrometer range, Van
der Waals’ forces, electrostatic interactions and hydrogen bonding can cause two materials to adhere to one another. A force is
required to overcome these interactions and cause the surfaces to move across one another.
Like kinetic friction, the force of static friction is given by a coefficient multiplied by the normal force. The normal force is the
force of the surface pushing up on the object, which is usually equal to the object’s weight. The coefficient of static friction is
usually greater than the coefficient of kinetic friction and is usually represented by μ .
s
Putting these elements together gives the maximum force of static friction as:
Fs = μ Fn (4.7.1)
s
As with all frictional forces, the force of friction can never exceed the force applied. Thus the force of static friction will vary
between 0 and F depending on the strength of the applied force. Any force smaller than F attempting to slide one surface over
s n s n
the other is opposed by a frictional force of equal magnitude and opposite direction. Any force larger than that overcomes the force
of static friction and causes sliding to occur. The instant sliding occurs, static friction is no longer applicable—the friction between
the two surfaces is then called kinetic friction.
Static Friction: To move a block at rest on a surface, a force must be applied which is great enough to overcome the force of static
friction.
4.7.4 https://phys.libretexts.org/@go/page/15229
learning objectives
Calculate the force of friction on an incline
As always, the frictional force resists motion. If the block is being pushed up the incline the friction force points down the incline.
If the block is being pulled down the incline, the friction force will hold the block up.
the force of gravity, the sum of the forces is still 0. The force of friction can never exceed the other forces acting on it. The
frictional forces only act to counter motion.
Drag
The drag force is the resistive force felt by objects moving through fluids and is proportional to the square of the object’s speed.
learning objectives
Relate the magnitude of drag force to the speed of an object
Another interesting force in everyday life is the force of drag on an object when it is moving in a fluid (either a gas or a liquid).
You feel the drag force when you move your hand through water. You might also feel it if you move your hand during a strong
wind. The faster you move your hand, the harder it is to move. You feel a smaller drag force when you tilt your hand so only the
side goes through the air—you have decreased the area of your hand that faces the direction of motion.
Like friction, the drag force always opposes the motion of an object. Unlike simple friction, the drag force is proportional to some
function of the velocity of the object in that fluid. This functionality is complicated and depends upon the shape of the object, its
size, its velocity, and the fluid it is in. Aerodynamic objects tend to have small surface areas and be designed to have low drag
coefficients.
For most large objects such as bicyclists, cars, and baseballs not moving too slowly, the magnitude of the drag force F is found to
D
be proportional to the square of the speed of the object. We can write this relationship mathematically as FD ∝ v . When taking
2
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into account other factors, this relationship becomes FD =
1
2C
, where C is known as the drag coefficient, a unit-less number
2
r hoAv
that represents the aerodynamic properties of the object, A is the cross-sectional area of the object which is facing the direction of
motion, and ρ is the density of the fluid the object is moving through.
Aerodynamic Shape: From racing cars to bobsled racers, aerodynamic shaping is crucial to achieving top speeds. Bobsleds are
designed for speed. They are shaped like a bullet with tapered fins. (credit: U.S. Army, via Wikimedia Commons)
Retarding and Drag Forces: A brief look at retarding (drag) forces in physics, for students in introductory physics classes that use
calculus. This video walks through a single scenario of an object experiencing a drag force where the drag force is proportional to
the object’s velocity.
4.7.6 https://phys.libretexts.org/@go/page/15229
Stress and Strain
The ratio of force to area F
A
is called stress and the ratio of change in length to length ΔL
L
is called the strain.
learning objectives
Explain how forces affects the shape of an object
We now move from consideration of forces that affect the motion of an object (such as friction and drag) to those that affect an
object’s shape. If a bulldozer pushes a car into a wall, the car will not move past the wall but it will noticeably change shape. A
change in shape due to the application of a force is a deformation. Even very small forces are known to cause some deformation.
For small deformations, two important characteristics are observed. First, the object returns to its original shape when the force is
removed—that is, the deformation is elastic for small deformations. Second, the size of the deformation is proportional to the force
—that is, for small deformations, Hooke’s law is obeyed. In equation form, Hooke’s law is given by F = k ⋅ ΔL where ΔL is the
change in length and kk is a constant which depends on the material properties of the object.
Deformations come in several types: changes in length (tension and compression), sideways shear (stress), and changes in volume.
Tension/Compression: Tension: The rod is stretched a length ΔL when a force is applied parallel to its length. (b) Compression:
The same rod is compressed by forces with the same magnitude in the opposite direction. For very small deformations and uniform
materials, ΔL is approximately the same for the same magnitude of tension or compression. For larger deformations, the cross-
sectional area changes as the rod is compressed or stretched.
The ratio of force to area F
A
is called stress and the ratio of change in length to length ΔL
L
is called the strain.
Stress and strain are related to each other by a constant called Young’s Modulus or the elastic modulus which varies depending on
the material. Using Young’s Modulus the relation between stress and strain is given by: stress = Y ⋅ strain .
A material with a high elastic modulus is said to have high tensile strength. Such materials are very resistant to being stretched and
require a large amount of force to deform a small amount.
Translational Equilibrium
An object is said to be in equilibrium when there is no external net force acting on it.
learning objectives
Assess the role each type of equilibrium plays in mechanical devices
4.7.7 https://phys.libretexts.org/@go/page/15229
We are surrounded by great engineering architectures and mechanical devices, which are at rest in the frame of reference of Earth.
A large part of engineering creations are static objects. Yet we also seek equilibrium of moving objects like that of floating ship,
airplane cruising at high speed, and such other moving mechanical devices. In both cases – static or dynamic – net external forces
and torques are zero.
A body is said to be in mechanical equilibrium when net external force is equal to zero and net external torque is also zero.
Mathematically,
→ →
Σ F ext = 0 and Σ τ ext = 0 (4.7.3)
Since there is no net force on the object, the object does not accelerate. This implies two types of possible equilibrium. The first
type, where all particles in the system are at rest and do have not velocity, is known as static equilibrium. In the second type, the
object has a velocity, but since there are no net forces acting on it, the velocity remains constant. In the second case, the particle is
said to be in dynamic equilibrium. Static or dynamic, these kinds of equilibrium can be categorized as translational equilibrium.
Examples of translational equilibrium are all around us. A book resting on a table is pushing down on the table with the force of its
weight. The table, in turn, is pushing back on the book, keeping the book from falling through the table. Since neither the table nor
the book are moving, this is an example of static equilibrium. The force of gravity on the book is perfectly counteracted by the
force of the table pushing on it.
Forces Acting on an Object at Rest: A force diagram showing the forces acting on an object at rest on a surface. Notice that the
amount of force that the table is pushing upward on the object (the N vector) is equal to the downward force of the object’s weight
(shown here as mg, as weight is equal to the object’s mass multiplied by the acceleration due to gravity): because these forces are
equal, the object is in a state of equilibrium (all the forces acting on it balance to zero).
An example of dynamic (or mechanical) equilibrium is an object sliding down a wedge. The force of gravity pulls the object down
the wedge, but it is counteracted by the force of friction between the wedge and the object. If the force of friction is equal to the
force of gravity, the object will proceed at a constant velocity.
Forces on a Table: These six forces are in equilibrium. The four forces of the table leg counteract the force of the table and the
object pushing on them.
Connected Objects
Forces can be transferred from one object to another through connections.
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learning objectives
Analyze the affect a rigid connection has on the movement of objects
The physics of connected objects is very similar to physics of simple objects. There are a variety of ways objects can be connected
to each other, and a corresponding variety of mathematical ways to model such connections.
The simplest form of connection is a perfectly rigid connection. If two objects are connected by a perfectly rigid connector then
they may be thought of as the same object. Perfectly rigid connectors cannot stretch nor deform, and transfer forces instantaneously
from one side of the connection to the other. For example, given two blocks (both of mass 1 kg) connected by a perfectly rigid bar,
if the first block is pulled with a force of 1 Newton, then both blocks will accelerate at the same time and the same acceleration. In
this case the acceleration is m/s —the same as if a force of mass 2 kg is exerted on one object. Thus it can be said that a
1
2
2
perfectly rigid connection makes two objects into one large object. Of course, perfectly rigid connections do not exist in nature.
Some deformation will always exist in any object as force travels along it. However, many materials are sufficiently rigid, so that
using the perfectly rigid approximation is useful for simplicity’s sake.
One can think of the force transferring through the connection by means of the “tension” force. Tension is the pulling force exerted
by a string, chain, or similar connector on another object. If two objects are connected by a string, a force exerted on one is
balanced by a tension force in the string which pulls on the other. Of course, if the tension force is greater than the rope can
withstand, the rope will break.
Tension Forces: The forces involved in supporting a ball by a rope. Tension is the force of the rope on the scaffold, the force of the
rope on the ball, and the balanced forces acting on and produced by segments of the rope.
Circular Motion
An object in circular motion undergoes acceleration due to centripetal force in the direction of the center of rotation.
learning objectives
Develop an understanding of uniform circular motion as an indicator for net external force
Uniform circular motion describes the motion of an object along a circle or a circular arc at constant speed. It is the basic form of
rotational motion in the same way that uniform linear motion is the basic form of translational motion. However, the two types of
motion are different with respect to the force required to maintain the motion.
Let us consider Newton’s first law of motion. It states that an object will maintain a constant velocity unless a net external force is
applied. Therefore, uniform linear motion indicates the absence of a net external force. On the other hand, uniform circular motion
requires that the velocity vector of an object constantly change direction. Since the velocity vector of the object is changing, an
acceleration is occurring. Therefore, uniform circular motion indicates the presence of a net external force.
In uniform circular motion, the force is always perpendicular to the direction of the velocity. Since the direction of the velocity is
continuously changing, the direction of the force must be as well.
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The direction of the velocity along the circular trajectory is tangential. The perpendicular direction to the circular trajectory is,
therefore, the radial direction. Therefore, the force (and therefore the acceleration) in uniform direction motion is in the radial
direction. For this reason, acceleration in uniform circular motion is recognized to “seek the center” — i.e., centripetal force.
The equation for the acceleration aa required to sustain uniform circular motion is:
2
v
a = (4.7.4)
r
where mm is the mass of the object, v is the velocity of the object, and r is the radius of the circle. Consequently, the net external
force F required to sustain circular motion is:
net
2
m⋅v
Fnet = (4.7.5)
r
Uniform Circular Motion: In uniform circular motion, the centripetal force is perpendicular to the velocity. The centripetal force
points toward the center of the circle, keeping the object on the circular track.
Key Points
Friction is the force that resists relative motion between two surfaces sliding across each other. Friction converts kinetic energy
into heat.
Drag force is the force that resists motion of an object traveling through a fluid such as air or water. Drag force is proportional
to the velocity of the object traveling.
Deformation forces are forces caused by stretching or compressing a material. Some examples would be springs or elastics.
Kinetic (or dynamic) friction occurs when two objects are moving relative to each other and rub together (like a sled on the
ground).
The force of friction can be represented by an equation F f riction = μF where μ is the coefficient of friction and is a unitless
n
normal force.
Motion on an incline is resisted by friction.
The frictional force on an incline is dependent on the angle of the incline. F = μmg cos (θ) is the maximum friction force on
f
an incline.
If the friction force is greater than or equal to the forces in the direction of motion, then the net force is 0 and the object is in
equilibrium.
Objects moving through a fluid feel a force which resists their motion. This force is known as the drag force.
The drag force is proportional to the square of the velocity of the object relative to the fluid.
The equation for drag is F = C Av . C is a constant called the drag coefficient. ρρ is the density of the fluid. A is the
D
1
2
ρ
2
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The ratio of force to area is called stress and the ratio of change in length to length
F
A
ΔL
L
is called the strain.
Stress and strain are related to each other by a constant called Young’s Modulus or the elastic modulus which varies depending
on the material. Using Young’s Modulus the relation between stress and strain is given by: stress = Y ⋅ strain .
A material with a high elastic modulus is said to have high tensile strength. Such materials are very resistant to being stretched
and require a large amount of force to deform a small amount.
When there is no external net force on an object, the object is said to be in equilibrium.
When an object is in equilibrium, it does not accelerate. If it had a velocity, the velocity remains constant; if it was at rest, it
remains at rest.
An equilibrium in motion is known as dynamic equilibrium; an equilibrium at rest is a static equilibrium.
If two objects are connected, a force on one has an effect on the other.
Connections can often be approximated as completely rigid. In completely rigid cases, the connection does not deform and the
force is transferred instantaneously.
Tension is the force of a rope or cable or other connector on the object it is connected to. It is one way force is transferred
between objects.
An object that is undergoing circular motion has a velocity vector that is constantly changing direction.
The force that is needed to maintain circular motion points toward the center of the circular path. It is therefore known as the
centripetal force.
The velocity of an object in circular motion is always tangent to the circle, and the centripetal force is always perpendicular to
the velocity.
Key Terms
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
static: Fixed in place; having no motion.
kinetic: Of or relating to motion
friction: A force that resists the relative motion or tendency to such motion of two bodies in contact.
incline: A slope.
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
fluid: Any substance which can flow with relative ease, tends to assume the shape of its container, and obeys Bernoulli’s
principle; a liquid, gas or plasma.
strain: The amount by which a material deforms under stress or force, given as a ratio of the deformation to the initial
dimension of the material and typically symbolized by ϵϵ is termed the engineering strain. The true strain is defined as the
natural logarithm of the ratio of the final dimension to the initial dimension.
stress: The internal distribution of force per unit area (pressure) within a body reacting to applied forces which causes strain or
deformation and is typically symbolized by σ.
dynamic: Changing; active; in motion.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
rigid: Stiff, rather than flexible.
tangent: a straight line touching a curve at a single point without crossing it at that point
perpendicular: at or forming a right angle (to).
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CHAPTER OVERVIEW
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1
5.1: Introduction to UCM and Gravitation
Kinematics of UCM
Uniform circular motion is a motion in a circular path at constant speed.
learning objectives
Relate centripetal force and centripetal acceleration to uniform circular motion
Angular Quantities
Under uniform circular motion, angular and linear quantities have simple relations. When objects rotate about some axis, each
point in the object follows a circular arc. The rotation angle is the amount of rotation and is analogous to linear distance. We define
the rotation angle Δθ to be the ratio of the arc length to the radius of curvature:
Δs
Δθ = (5.1.1)
r
Angle θθ and Arc Length ss: The radius of a circle is rotated through an angle Δθ. The arc length Δs is described on the
circumference.
We define angular velocity ωω as the rate of change of an angle. In symbols, this is ω =
Δθ
Δt
, where an angular rotation Δθ takes
place in a time Δt. From the relation of s and (Δs = rΔθ ), we see:
Δs Δθ
v = =r = rω (5.1.2)
Δt Δt
Under uniform circular motion, the angular velocity is constant. The acceleration can be written as:
2
dv dr 2
v
ac = =ω = ωv = rω = (5.1.3)
dt dt r
This acceleration, responsible for the uniform circular motion, is called centripetal acceleration.
Centripetal Force
Any force or combination of forces can cause a centripetal or radial acceleration. Just a few examples are the tension in the rope on
a tether ball, the force of Earth’s gravity on the Moon, friction between roller skates and a rink floor, a banked roadway’s force on a
car, and forces on the tube of a spinning centrifuge.
Any net force causing uniform circular motion is called a centripetal force. The direction of a centripetal force is toward the center
of curvature, the same as the direction of centripetal acceleration. According to Newton’s second law of motion, net force is mass
times acceleration. For uniform circular motion, the acceleration is the centripetal acceleration: a = a . Thus, the magnitude of
c
2
v
2
Fc = mac = m = mrω (5.1.4)
r
5.1.1 https://phys.libretexts.org/@go/page/14459
Dynamics of UCM
Newton’s universal law of gravitation states that every particle attracts every other particle with a force along a line joining them.
learning objectives
Relate Kepler’s laws to Newton’s universal law of gravitation
The gravitational force is responsible for artificial satellites orbiting the Earth. The Moon’s orbit about Earth, the orbits of planets,
asteroids, meteors, and comets about the Sun are other examples of gravitational orbits. Historically, Kepler discovered his 3 laws
(called Kepler’s law of planetary motion) long before the days of Newton. Kepler devised his laws after careful study (over some
20 years) of a large amount of meticulously recorded observations of planetary motion done by Tycho Brahe (1546–1601).
Kepler’s Laws
The orbit of each planet about the Sun is an ellipse with the Sun at one focus.
Each planet moves so that an imaginary line drawn from the Sun to the planet sweeps out equal areas in equal times.
The ratio of the squares of the periods of any two planets about the Sun is equal to the ratio of the cubes of their average
distances from the Sun.
Kepler’s Second Law: The shaded regions have equal areas. It takes equal times for mm to go from A to B , from C to D, and
from E to F . The mass mm moves fastest when it is closest to M . Kepler’s second law was originally devised for planets orbiting
the Sun, but it has broader validity.
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Ellipses and Kepler’s First Law: (a) An ellipse is a closed curve such that the sum of the distances from a point on the curve to
the two foci (f and f ) is a constant. You can draw an ellipse as shown by putting a pin at each focus, and then placing a string
1 2
around a pencil and the pins and tracing a line on paper. A circle is a special case of an ellipse in which the two foci coincide (thus
any point on the circle is the same distance from the center). (b) For any closed gravitational orbit, m follows an elliptical path with
M at one focus. Kepler’s first law states this fact for planets orbiting the Sun.
T is the period (time for one orbit) and r is the average radius. We shall derive Kepler’s third law, starting with Newton’s laws of
motion and his universal law of gravitation. We will assume a circular path (not an elliptical one) for simplicity.
Let us consider a circular orbit of a small mass mm around a large mass M , satisfying the two conditions stated at the beginning of
this section. Gravity supplies the centripetal force to mass m. Therefore, for a uniform circular motion:
2
mM v
G = mac = m (5.1.7)
2
r r
Now, to get at Kepler’s third law, we must get the period T into the equation. By definition, period T is the time for one complete
orbit. Now the average speed v is the circumference divided by the period:
2πr
v = (5.1.9)
T
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Substituting this into the previous equation gives:
2 2
M 4π r
G = (5.1.10)
2
r T
2
2
4π 3
T = r (5.1.11)
GM
learning objectives
Derive θ for an ideally banked curve for speed
Overview
As an example of a uniform circular motion and its application, let us now consider banked curves, where the slope of the road
helps you negotiate the curve. The greater the angle θ , the faster you can take the curve. Race tracks for bikes as well as cars, for
example, often have steeply banked curves. In an “ideally banked curve,” the angle θ is such that you can negotiate the curve at a
certain speed without the aid of friction between the tires and the road. We will derive an expression for θ for an ideally banked
curve for speed v and consider an example related to it.
Car on a Banked Curve: The car on this banked curve is moving away and turning to the left.
Above is a free body diagram for a car on a frictionless banked curve. The only two external forces acting on the car are its weight
ww and the normal force of the road N . (A frictionless surface can only exert a force perpendicular to the surface—that is, a
normal force. ) These two forces must add to give a net external force that is horizontal toward the center of curvature and has
2
magnitude mv
r
. Only the normal force has a horizontal component, and so this must equal the centripetal force—that is:
5.1.4 https://phys.libretexts.org/@go/page/14459
2
mv
N sin θ = (5.1.12)
r
Because the car does not leave the surface of the road, the net vertical force must be zero, meaning that the vertical components of
the two external forces must be equal in magnitude and opposite in direction. From the figure, we see that the vertical component
of the normal force is N cos θ, and the only other vertical force is the car’s weight. These must be equal in magnitude, thus:
N cos θ = mg (5.1.13)
Key Points
Under uniform circular motion, angular and linear quantities have simple relations. The length of an arc is proportional to the
rotation angle and the radius. Also, v = rω .
2
The acceleration responsible for the uniform circular motion is called centripetal acceleration. It is given as a = rω =
c
2 v
r
.
Any net force causing uniform circular motion is called a centripetal force. The direction of a centripetal force is toward the
2
r
2
The gravitational force is directly proportional to the product of their masses and inversely proportional to the square of the
distance between them.
Kepler discovered laws describing planetary motion long before the days of Newton, purely based on the observations of Tycho
Brahe.
Kepler’s laws can be derived from the Newton’s universal law of gravitation and his equation of motion.
For ideal banking, the net external force equals the horizontal centripetal force in the absence of friction.
For ideal banking, the components of the normal force NN in the horizontal and vertical directions must equal the centripetal
force and the weight of the car, respectively.
2
rg
Key Terms
centripetal: Directed or moving towards a center.
asteroid: A naturally occurring solid object, which is smaller than a planet and is not a comet, that orbits a star.
planet: A large body which directly orbits any star (or star cluster) but which has not attained nuclear fusion.
normal force: Any force acting normal, to a surface, or perpendicular to the tangent plane.
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5.2: Non-Uniform Circular Motion
Overview of Non-Uniform Circular Motion
Non-uniform circular motion denotes a change in the speed of a particle moving along a circular path.
learning objectives
Explain when a particle undergoes non-uniform circular motion
What do we mean by non-uniform circular motion? The answer lies in the definition of uniform circular motion, which is a circular
motion with constant speed. It follows then that non-uniform circular motion denotes a change in the speed of the particle moving
along the circular path. Note especially the change in the velocity vector sizes, denoting change in the magnitude of velocity.
Diagram of non-uniform circular motion: In non-uniform circular motion, the magnitude of the angular velocity changes over
time.
The change in direction is accounted by radial acceleration ( centripetal acceleration ), which is given by following relation:
2
a =
r . The change in speed has implications for radial (centripetal) acceleration. There are two possibilities:
v
1: The radius of circle is constant (like in the motion along a circular rail or motor track). A change in v will change the magnitude
of radial acceleration. This means that the centripetal acceleration is not constant, as is the case with uniform circular motion. The
greater the speed, the greater the radial acceleration. A particle moving at higher speed will need a greater radial force to change
direction and vice-versa when the radius of the circular path is constant.
2: The radial (centripetal) force is constant (like a satellite rotating about the earth under the influence of a constant force of
gravity). The circular motion adjusts its radius in response to changes in speed. This means that the radius of the circular path is
variable, unlike the case of uniform circular motion. In any eventuality, the equation of centripetal acceleration in terms of “speed”
and “radius” must be satisfied. The important thing to note here is that, although change in speed of the particle affects radial
acceleration, the change in speed is not affected by radial or centripetal force. We need a tangential force to affect the change in the
magnitude of a tangential velocity. The corresponding acceleration is called tangential acceleration.
In either case, the angular velocity in non-uniform circular motion is not constant as ω = v
r
and v is varying.
Key Points
In non- uniform circular motion, the size of the velocity vector (speed) changes, denoting change in the magnitude of velocity.
The change in speed has implications for radial ( centripetal ) acceleration. There are two possibilities: 1) the radius of the circle
is constant; or 2) the radial (centripetal) force is constant.
In either case, the angular velocity in non-uniform circular motion is not constant, as ω = , and v varies.
v
Key Terms
radial: Moving along a radius.
centripetal: Directed or moving towards a center.
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5.3: Velocity, Acceleration, and Force
Rotational Angle and Angular Velocity
The rotational angle is a measure of how far an object rotates, and angular velocity measures how fast it rotates.
learning objectives
Express the relationship between the rotational angle and the distance
R
, it is usually more convenient to use radians.
Angle θ and Arc Length s: The radius of a circle is rotated through an angle Δθ. The arc length Δs is described on the
circumference.
The speed at which the object rotates is given by the angular velocity, which is the rate of change of the rotational angle with
respect to time. Although the angle itself is not a vector quantity, the angular velocity is a vector. The direction of the angular
velocity vector is perpendicular to the plane of rotation, in a direction which is usually specified by the right-hand rule. Angular
acceleration gives the rate of change of angular velocity. The angle, angular velocity, and angular acceleration are very useful in
describing the rotational motion of an object.
The Direction of Angular Velocity: The angular velocity describes the speed of rotation and the orientation of the instantaneous
axis about which the rotation occurs. The direction of the angular velocity will be along the axis of rotation. In this case (counter-
clockwise rotation), the vector points upwards.
When the axis of rotation is perpendicular to the position vector, the angular velocity may be calculated by taking the linear
velocity v of a point on the edge of the rotating object and dividing by the radius. This will give the angular velocity, usually
denoted by ω, in terms of radians per second.
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Angular Velocity: A fly on the edge of a rotating object records a constant velocity v . The object is rotating with an angular
velocity equal to . v
Centripetial Acceleration
Centripetal acceleration is the constant change in velocity necessary for an object to maintain a circular path.
learning objectives
Express the centripetal acceleration in terms of rotational velocity
Overview
As mentioned in previous sections on kinematics, any change in velocity is given by an acceleration. Often the changes in velocity
are changes in magnitude. When an object speeds up or slows down this is a change in the objects velocity. Changes in the
magnitude of the velocity match our intuitive and every day usage of the term accelerate. However, because velocity is a vector, it
also has a direction. Therefore, any change in the direction of travel of an object must also be met with an acceleration.
Uniform circular motion involves an object traveling a circular path at constant speed. Since the speed is constant, one would not
usually think that the object is accelerating. However, the direction is constantly changing as the object traverses the circle. Thus, it
is said to be accelerating. One can feel this acceleration when one is on a roller coaster. Even if the speed is constant, a quick turn
will provoke a feeling of force on the rider. This feeling is an acceleration.
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High School Physics - Centripetal Acceleration
Centripetal Acceleration: A brief overview of centripetal acceleration for high school physics students.
where v is the linear velocity of the object and r is the radius of the circle.
The centripetal acceleration may also be expressed in terms of rotational velocity as follows:
2
ac = ω r (5.3.2)
r
.
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Centripetal Acceleration: As an object moves around a circle, the direction of the velocity vector constantly changes.
Centripetal Force
A force which causes motion in a curved path is called a centripetal force (uniform circular motion is an example of centripetal
force).
learning objectives
Express the equations for the centripetal force and acceleration
A force that causes motion in a curved path is called a centripetal force. Uniform circular motion is an example of centripetal force
in action. It can be seen in the orbit of satellites around the earth, the tension in a rope in a game of tether ball, a roller coaster loop
de loop, or in a bucket swung around the body.
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High School Physics - Centripetal Force
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Centripetal force: As an object travels around a circular path at a constant speed, it experiences a centripetal force accelerating it
toward the center.
The equation for centripetal force is as follows:
2
mv
Fc = (5.3.3)
r
where: F is centripetal force, m is mass, v is velocity, and r is the radius of the path of motion.
c
From Newton’s second law F = m ⋅ a , we can see that centripetal acceleration is:
2
v
ac = (5.3.4)
r
Centripetal force can also be expressed in terms of angular velocity. Angular velocity is the measure of how fast an object is
traversing the circular path. As the object travels its path, it sweeps out an arc that can be measured in degrees or radians. The
equation for centripetal force using angular velocity is:
2
Fc = mrω (5.3.5)
Key Points
When an object rotates about an axis, the points on the edge of the object travel in arcs.
The angle these arcs sweep out is called the rotational angle, and it is usually represented by the symbol theta.
A measure of how quickly the object is rotating, with respect to time, is called the angular velocity. It is usually represented by a
Greek omega symbol. Like its counterpart linear velocity, it is a vector.
For an object to maintain circular motion it must constantly change direction.
Since velocity is a vector, changes in direction constitute changes in velocity.
A change in velocity is known as an acceleration. The change in velocity due to circular motion is known as centripetal
acceleration.
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Centripetal acceleration can be calculated by taking the linear velocity squared divided by the radius of the circle the object is
traveling along.
When an object is in uniform circular motion, it is constantly changing direction, and therefore accelerating. This is angular
acceleration.
A force acting on the object in uniform circular motion (called centripetal force) is acting on the object from the center of the
circle.
Key Terms
radians: The angle subtended at the centre of a circle by an arc of the circle of the same length as the circle’s radius.
acceleration: The amount by which a speed or velocity increases (and so a scalar quantity or a vector quantity).
circular motion: Motion in such a way that the path taken is that of a circle.
velocity: A vector quantity that denotes the rate of change of position with respect to time, or a speed with a directional
component.
centripetal: Directed or moving towards a center.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
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5.4: Types of Forces in Nature
Tides
Tides are the rise and fall of sea levels due to the effects of the gravity exerted by the moon and the sun, and the rotation of the
Earth.
learning objectives
Explain factors that influence the times and amplitude of the tides at a locale
Tides are the rise and fall of sea levels due to the effects of gravitational forces exerted by the moon and the sun when combined
with the rotation of the Earth. Tides occur to varying degrees and frequency, depending on location. Shorelines where two almost
equally high tides and two low tides occur each day experience a semi- diurnal tide. The occurrence of only one high and one low
tide each day is known as diurnal tide. A mixed tide referes to the daily occurrence of two uneven tides, or perhaps one high and
one low tide. The times and amplitude of tides at various locales are influenced by the alignment of the sun and moon, the pattern
of tides in the deep ocean, the shape of the coastline, and other forces.
Earth’s tides.: Schematic of the lunar portion of earth’s tides showing (exaggerated) high tides at the sublunar and antipodal points
for the hypothetical case of an ocean of constant depth with no land. There would also be smaller, superimposed bulges on the sides
facing toward and away from the sun.
Tidal Force
If we want to know the acceleration “felt” by an observer living on Earth due to the moon, a tricky part is that the Earth is not an
inertial frame of reference because it is in “free fall” with respect to the moon. Given this, in order to figure out the force observed,
we must subtract the acceleration of the (Earth) frame itself. The tidal force produced by the moon on a small particle located on
Earth is the vector difference between the gravitational force exerted by the moon on the particle, and the gravitational force that
would be exerted if it were located at the Earth’s center of mass.
5.4.1 https://phys.libretexts.org/@go/page/14462
Moon’s Gravity on the Earth: Top picture shows the gravity force due to the Moon at different locations F on Earth. Bottom
r
As diagramed below, this is equivalent to subtracting the “red” vector from the “black” vectors on the surface of the Earth in the
top picture, leading to the “differential” force represented by the bottom picture. Thus, the tidal force depends not on the strength of
the lunar gravitational field, but on its gradient (which falls off approximately as the inverse cube of the distance to the originating
gravitational body).
On average, the solar gravitational force on the Earth is 179 times stronger than the lunar, but because the sun is on average 389
times farther from the Earth its field gradient is weaker. The solar tidal force is 46% as large as the lunar. More precisely, the lunar
tidal acceleration (along the moon-Earth axis, at the Earth’s surface) is about 1.1 ⋅ 10 g, while the solar tidal acceleration (along
−7
the sun-Earth axis, at the Earth’s surface) is about 0.52 ⋅ 10 g, where g is the gravitational acceleration at the Earth’s surface.
−7
Venus has the largest effect of the other planets, at 0.000113 times the solar effect.
Tidal Energy
Energy of tides can be extracted by two means: inserting a water turbine into a tidal current, or building ponds that release/admit
water through a turbine. In the first case, the energy amount is entirely determined by the timing and tidal current magnitude, but
the best currents may be unavailable because the turbines would obstruct ships. In the second case, impoundment dams are
expensive to construct, natural water cycles are completely disrupted, as is ship navigation. However, with multiple ponds, power
can be generated at chosen times. Presently, there are few installed systems for tidal power generation (most famously, La Rance by
Saint Malo, France), as many difficulties are involved. Aside from environmental issues, simply withstanding corrosion and
biological fouling pose engineering challenges.
Tidal Energy Generator: Tidal energy generator that works like a wind turbine, but with the ocean currents providing the energy.
The circle in the middle is the turbine. The contraption travels up and down the two legs just like a lift and sits on the sea floor
when in use.
Unlike with wind power systems, tidal power proponents point out that generation levels can be reliably predicted (save for
weather effects). While some generation is possible for most of the tidal cycle, in practice, turbines lose efficiency at lower
operating rates. Since the power available from a flow is proportional to the cube of the flow speed, the times during which high
power generation is possible are brief.
learning objectives
Formulate relationship between the Coriolis force, mass of an object, and speed in the rotating frame
The Coriolis effect is a deflection of moving objects when they are viewed in a rotating reference frame. In a reference frame with
clockwise rotation, the deflection is to the left of the motion of the object; in one with counter-clockwise rotation, the deflection is
to the right. Although recognized previously by others, the mathematical expression for the Coriolis force appeared in an 1835
paper by French scientist Gaspard-Gustave Coriolis, in connection with the theory of water wheels. Early in the 20th century, the
term “Coriolis force” began to be used in connection with meteorology.
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Frames of Reference: In the inertial frame of reference (upper part of the picture), the black object moves in a straight line.
However, the observer (red dot) who is standing in the rotating/non-inertial frame of reference (lower part of the picture) sees the
object as following a curved path due to the Coriolis and centrifugal forces present in this frame.
Newton’s laws of motion govern the motion of an object in a (non-accelerating) inertial frame of reference. When Newton’s laws
are transformed to a uniformly rotating frame of reference, the Coriolis and centrifugal forces appear. Both forces are proportional
to the mass of the object. The Coriolis force is proportional to the rotation rate, and the centrifugal force is proportional to its
square. The Coriolis force acts in a direction perpendicular to the rotation axis and to the velocity of the body in the rotating frame.
It is proportional to the object’s speed in the rotating frame. These additional forces are termed inertial forces, fictitious forces, or
pseudo-forces. They allow the application of Newton’s laws to a rotating system. They are correction factors that do not exist in a
non-accelerating or inertial reference frame.
Perhaps the most commonly encountered rotating reference frame is the Earth. The Coriolis effect is caused by the rotation of the
Earth and the inertia of the mass experiencing the effect. Because the Earth completes only one rotation per day, the Coriolis force
is quite small. Its effects generally become noticeable only for motions occurring over large distances and long periods of time,
such as large-scale movements of air in the atmosphere or water in the ocean. Such motions are constrained by the surface of the
earth, so generally only the horizontal component of the Coriolis force is important. This force causes moving objects on the
surface of the Earth to be deflected in a clockwise sense (with respect to the direction of travel) in the northern hemisphere and in a
counter-clockwise sense in the southern hemisphere. Rather than flowing directly from areas of high pressure to low pressure, as
they would in a non-rotating system, winds and currents tend to flow to the right of this direction north of the equator and to the left
of this direction south of it. This effect is responsible for the rotation of large cyclones.
Coriolis Force: This low-pressure system over Iceland spins counter-clockwise due to balance between the Coriolis force and the
pressure gradient force.
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Other Geophysical Applications
Tidal and Coriolis forces may not be obvious over a small time-space scale, but they are important in meteorology, navigation, and
fishing.
learning objectives
Identify fields that have to take into account the tidal and Coriolis forces
We have studied tidal and Coriolis forces previously. To review, the tidal force is responsible for the tides — it is a “differential
force,” due to a secondary effect of the force of gravity. The Coriolis force is a fictitious force, representing a deflection of moving
objects when they are viewed in a rotating reference frame of the Earth. Although their effects may not be obvious over a small
time-space scale, these forces are important in such contexts as meteorology, navigation, fishing, and others.
The Tides
Tidal flows are important for marine navigation, and significant errors in position occur if they are not accounted for. Tidal heights
are also important; for example, many rivers and harbors have a shallow “bar” at the entrance to prevent boats with significant draft
from entering at low tide. Until the advent of automated navigation, competence in calculating tidal effects was important to naval
officers. The certificate of examination for lieutenants in the Royal Navy once declared that the prospective officer was able to
“shift his tides. ”
Tidal Indicator: Tidal Indicator, Delaware River, Delaware c. 1897. In the moment pictured, the tide is 1.25 feet above mean low
water and is still falling, as indicated by the pointing of the arrow. The indicator is powered by a system of pulleys, cables, and a
float
5.4.4 https://phys.libretexts.org/@go/page/14462
Flow Representation: A schematic representation of flow around a low-pressure area in the Northern Hemisphere. The pressure-
gradient force is represented by blue arrows and the Coriolis acceleration (always perpendicular to the velocity) by red arrows
Key Points
The tidal force depends not on the strength of the lunar gravitational field itself, but on its gradient, which falls off
approximately as the inverse cube of the distance to the originating gravitational body. This is because the tidal force felt by an
observer on Earth is a differential force.
The times and amplitude of the tides at a locale are influenced by several factors, such as the alignment of the sun and moon,
the pattern of tides in the deep ocean, the shape of the coastline, and others forces.
Energy of tides can be extracted by two means: inserting a water turbine into a tidal current, or building ponds that
release/admit water through a turbine.
When Newton’s laws are transformed to a uniformly rotating frame of reference, the Coriolis and centrifugal forces appear.
The Coriolis force acts in a direction perpendicular to the rotation axis and to the velocity of the body in the rotating frame; it is
proportional to the object’s mass and speed in the rotating frame.
The Coriolis effect is caused by the rotation of the Earth and the inertia of the mass experiencing the effect.
Tidal flows are important for marine navigation, and significant errors in position occur if they are not accounted for.
The Coriolis force is quite small, and its effects generally become noticeable only when we are dealing with motions occurring
over large distances and long periods of time, such as large-scale movements of air in the atmosphere or water in the ocean.
The tidal force is responsible for the tides. It is a “differential force,” due to a secondary effect of the force of gravity. The
Coriolis force is a fictitious force, representing a deflection of moving objects when they are viewed in a rotating reference
frame of the Earth.
Key Terms
inertial frame: A frame of reference that describes time and space homogeneously, isotropically, and in a time-independent
manner.
diurnal: Having a daily cycle that is completed every 24 hours, usually referring to tasks, processes, tides, or sunrise to sunset.
gradient: The rate at which a physical quantity increases or decreases relative to change in a given variable, especially distance.
fictitious force: an apparent force that acts on all masses in a non-inertial frame of reference, such as a rotating reference frame
centrifugal force: the apparent outward force that draws a rotating body away from the center of rotation
ballistics: the science of mechanics that deals with the flight, behavior, and effects of projectiles, especially bullets, gravity
bombs, rockets, or the like
meteorology: the interdisciplinary scientific study of the atmosphere
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5.5: Newton’s Law of Universal Gravitation
The Law of Universal Gravitation
Objects with mass feel an attractive force that is proportional to their masses and inversely proportional to the square of the
distance.
learning objectives
Express the Law of Universal Gravitation in mathematical form
While an apple might not have struck Sir Isaac Newton’s head as myth suggests, the falling of one did inspire Newton to one of the
great discoveries in mechanics: The Law of Universal Gravitation. Pondering why the apple never drops sideways or upwards or
any other direction except perpendicular to the ground, Newton realized that the Earth itself must be responsible for the apple’s
downward motion.
Theorizing that this force must be proportional to the masses of the two objects involved, and using previous intuition about the
inverse-square relationship of the force between the earth and the moon, Newton was able to formulate a general physical law by
induction.
The Law of Universal Gravitation states that every point mass attracts every other point mass in the universe by a force pointing in
a straight line between the centers-of-mass of both points, and this force is proportional to the masses of the objects and inversely
proportional to their separation This attractive force always points inward, from one point to the other. The Law applies to all
objects with masses, big or small. Two big objects can be considered as point-like masses, if the distance between them is very
large compared to their sizes or if they are spherically symmetric. For these cases the mass of each object can be represented as a
point mass located at its center-of-mass.
While Newton was able to articulate his Law of Universal Gravitation and verify it experimentally, he could only calculate the
relative gravitational force in comparison to another force. It wasn’t until Henry Cavendish’s verification of the gravitational
constant that the Law of Universal Gravitation received its final algebraic form:
Mm
F =G (5.5.1)
r2
where F represents the force in Newtons, M and m represent the two masses in kilograms, and r represents the separation in
meters. G represents the gravitational constant, which has a value of 6.674 ⋅ 10 N(m/kg) . Because of the magnitude of G,
−11 2
Forces on two masses: All masses are attracted to each other. The force is proportional to the masses and inversely proportional to
the square of the distance.
learning objectives
Formulate the Shell Theorem for spherically symmetric objects
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However, most objects are not point particles. Finding the gravitational force between three-dimensional objects requires treating
them as points in space. For highly symmetric shapes such as spheres or spherical shells, finding this point is simple.
Diagram used in the proof of the Shell Theorem: This diagram outlines the geometry considered when proving The Shell
Theorem. In particular, in this case a spherical shell of mass M (left side of figure) exerts a force on mass m (right side of the
figure) outside of it. The surface area of a thin slice of the sphere is shown in color. (Note: The proof of the theorem is not
presented here. Interested readers can explore further using the sources listed at the bottom of this article.)
5.5.2 https://phys.libretexts.org/@go/page/14463
where it can be shown that M <d =
4
3
3
πd ρ
(ρ is the mass density of the sphere and we are assuming that it does not depend on the radius. That is, the sphere’s mass is
uniformly distributed.)
Therefore, combining the above two equations we get:
4
F = πGmρd (5.5.4)
3
which shows that mass mm feels a force that is linearly proportional to its distance, dd, from the sphere’s center of mass.
As in the case of hollow spherical shells, the net gravitational force that a solid sphere of uniformly distributed mass M exerts on a
body outside of it, is the vector sum of the gravitational forces acted by each shell of the sphere on the outside object. The resulting
net gravitational force acts as if mass M is concentrated on a point at the center of the sphere, which is the center of mass, or COM
(Statement 1 of Shell Theorem). More generally, this result is true even if the mass M is not uniformly distributed, but its density
varies radially (as is the case for planets).
learning objectives
Describe how gravitational force is calculated for the bodies with spatial extent
Newton’s law of universal gravitation states that every point mass in the universe attracts every other point mass with a force that is
directly proportional to the product of their masses, and inversely proportional to the square of the distance between them.
In modern language, the law states the following: Every point mass attracts every single other point mass by a force pointing along
the line intersecting both points. The force is proportional to the product of the two masses and inversely proportional to the square
of the distance between them:
m1 m2
F =G (5.5.5)
2
r
where F is the force between the masses, G is the gravitational constant, m is the first mass,
1 m2 is the second mass and r is the
distance between the centers of the masses.
If the bodies in question have spatial extent (rather than being theoretical point masses), then the gravitational force between them
is calculated by summing the contributions of the notional point masses which constitute the bodies. In the limit, as the component
point masses become “infinitely small”, this entails integrating the force (in vector form, see below) over the extents of the two
bodies.
In this way it can be shown that an object with a spherically-symmetric distribution of mass exerts the same gravitational attraction
on external bodies as if all the object’s mass were concentrated at a point at its center.
For points inside a spherically-symmetric distribution of matter, Newton’s Shell theorem can be used to find the gravitational force.
The theorem tells us how different parts of the mass distribution affect the gravitational force measured at a point located a distance
r from the center of the mass distribution:
0
1. The portion of the mass that is located at radii r < r causes the same force at r as if all of the mass enclosed within a sphere
0 0
of radius r was concentrated at the center of the mass distribution (as noted above).
0
2. The portion of the mass that is located at radii r > r exerts no net gravitational force at the distance r from the center. That is,
0 0
the individual gravitational forces exerted by the elements of the sphere out there, on the point at r , cancel each other out.
0
As a consequence, for example, within a shell of uniform thickness and density there is no net gravitational acceleration anywhere
within the hollow sphere. Furthermore, inside a uniform sphere the gravity increases linearly with the distance from the center; the
increase due to the additional mass is 1.5 times the decrease due to the larger distance from the center. Thus, if a spherically
symmetric body has a uniform core and a uniform mantle with a density that is less than of that of the core, then the gravity
2
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initially decreases outwardly beyond the boundary, and if the sphere is large enough, further outward the gravity increases again,
and eventually it exceeds the gravity at the core/mantle boundary.
The gravity of the Earth may be highest at the core/mantle boundary, as shown in Figure 1:
Gravitational Field of Earth: Diagram of the gravitational field strength within the Earth.
Key Points
Sir Isaac Newton’s inspiration for the Law of Universal Gravitation was from the dropping of an apple from a tree.
Newton’s insight on the inverse-square property of gravitational force was from intuition about the motion of the earth and the
moon.
The mathematical formula for gravitational force is F = G Mm
r2
where G is the gravitational constant.
Since force is a vector quantity, the vector summation of all parts of the shell contribute to the net force, and this net force is the
equivalent of one force measurement taken from the sphere’s midpoint, or center of mass (COM).
The gravitational force on an object within a hollow spherical shell is zero.
The gravitational force on an object within a uniform spherical mass is linearly proportional to its distance from the sphere’s
center of mass (COM).
Newton’s law of universal gravitation states that every point mass in the universe attracts every other point mass with a force
that is directly proportional to the product of their masses and inversely proportional to the square of the distance between them.
The second step in calculating earth’s mass came with the development of Newton’s law of universal gravitation.
By equating Newton’s second law with his law of universal gravitation, and inputting for the acceleration a the experimentally
verified value of 9.8 , the mass of earth is calculated to be 5.96 × 10 kg , making the earth’s weight calculable given any
m
s2
24
gravitational field.
The gravity of the Earth may be highest at the core/mantle boundary
Key Terms
induction: Use inductive reasoning to generalize and interpret results from applying Newton’s Law of Gravitation.
inverse: Opposite in effect or nature or order.
center of mass: The center of mass (COM) is the unique point at the center of a distribution of mass in space that has the
property that the weighted position vectors relative to this point sum to zero.
point mass: A theoretical point with mass assigned to it.
weight: The force on an object due to the gravitational attraction between it and the Earth (or whatever astronomical object it is
primarily influenced by).
gravitational force: A very long-range, but relatively weak fundamental force of attraction that acts between all particles that
have mass; believed to be mediated by gravitons.
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5.6: Kepler’s Laws
Kepler’s First Law
Kepler’s first law is: The orbit of every planet is an ellipse with the Sun at one of the two foci.
learning objectives
Apply Kepler’s first law to describe planetary motion
Definition
The orbit of every planet is an ellipse with the Sun at one of the two foci.
An ellipse is a closed plane curve that resembles a stretched out circle. Note that the Sun is not at the center of the ellipse, but at
one of its foci. The other focal point, f , has no physical significance for the orbit. The center of an ellipse is the midpoint of the
2
line segment joining its focal points. A circle is a special case of an ellipse where both focal points coincide.
Ellipses and Kepler’s First Law: (a) An ellipse is a closed curve such that the sum of the distances from a point on the curve to
the two foci (f and f ) is a constant. You can draw an ellipse as shown by putting a pin at each focus, and then placing a string
1 2
around a pencil and the pins and tracing a line on paper. A circle is a special case of an ellipse in which the two foci coincide (thus
any point on the circle is the same distance from the center). (b) For any closed gravitational orbit, m follows an elliptical path with
M at one focus. Kepler’s first law states this fact for planets orbiting the Sun.
How stretched out an ellipse is from a perfect circle is known as its eccentricity: a parameter that can take any value greater than or
equal to 0 (a circle) and less than 1 (as the eccentricity tends to 1, the ellipse tends to a parabola). The eccentricities of the planets
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known to Kepler varied from 0.007 (Venus) to 0.2 (Mercury). Minor bodies such as comets an asteroids (discovered after Kepler’s
time) can have very large eccentricities. The dwarf planet Pluto, discovered in 1929, has an eccentricity of 0.25.
Symbolically, an ellipse can be represented in polar coordinates as:
p
r = (5.6.1)
1 + ϵ cos θ
where (r, θ) are the polar coordinates (from the focus) for the ellipse, p is the semi-latus rectum, and ϵ is the eccentricity of the
ellipse. For a planet orbiting the Sun, r is the distance from the Sun to the planet and θ is the angle between the planet’s current
position and its closest approach, with the Sun as the vertex.
Orbit As Ellipse: Heliocentric coordinate system (r, θ) for ellipse. Also shown are: semi-major axis a, semi-minor axis b and
semi-latus rectum p; center of ellipse and its two foci marked by large dots. For θ = 0° , r = r and for θ = 180° , r = r
min . max
The semi-major axis a is the arithmetic mean between r min and r max :
rmax − a = a − rmin (5.6.4)
p
a = . (5.6.5)
2
1 −ϵ
The semi-minor axis b is the geometric mean between r min and r max :
rmax b
= (5.6.6)
b rmin
p
b= − −−−− (5.6.7)
√ 1 − ϵ2
The semi-latus rectum p is the harmonic mean between r min and r max :
1 1 1 1
− = − (5.6.8)
rmin p p rmax
2
pa = rmax ⋅ rmin = b . (5.6.9)
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The area of the ellipse is
A = πab (5.6.11)
Understanding Kepler’s 3 Laws and Orbits: In this video you will be introduced to Kepler’s 3 laws and see how they are
relevant to orbiting objects.
learning objectives
Apply Kepler’s second law to describe planetary motion
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Definition
A line joining a planet and the Sun sweeps out equal areas during equal intervals of time.
Understanding Kepler’s 3 Laws and Orbits: In this video you will be introduced to Kepler’s 3 laws and see how they are
relevant to orbiting objects.
In a small time the planet sweeps out a small triangle having base line and height. The area of this triangle is given by:
1
dA = ⋅ r ⋅ rdθ (5.6.12)
2
Now as the first law states that the planet follows an ellipse, the planet is at different distances from the Sun at different parts in its
orbit. So the planet has to move faster when it is closer to the Sun so that it sweeps equal areas in equal times.
The total area enclosed by the elliptical orbit is:
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A = πab (5.6.14)
Where θ˙ = dθ
dt
is the angular velocity, (using Newton notation for differentiation), and n=
2π
P
is the mean motion of the planet
around the Sun.
See below for an illustration of this effect. The planet traverses the distance between A and B, C and D, and E and F in equal times.
When the planet is close to the Sun it has a larger velocity, making the base of the triangle larger, but the height of the triangle
smaller, than when the planet is far from the Sun. One can see that the planet will travel fastest at perihelion and slowest at
aphelion.
Kepler’s Second Law: The shaded regions have equal areas. It takes equal times for m to go from A to B, from C to D, and from E
to F. The mass m moves fastest when it is closest to M. Kepler’s second law was originally devised for planets orbiting the Sun, but
it has broader validity.
learning objectives
Apply Kepler’s third law to describe planetary motion
Definition
The square of the orbital period of a planet is directly proportional to the cube of the semi-major axis of its orbit.
The third law, published by Kepler in 1619, captures the relationship between the distance of planets from the Sun, and their orbital
periods. Symbolically, the law can be expressed as
2 3
P ∝a , (5.6.16)
where P is the orbital period of the planet and a is the semi-major axis of the orbit (see ).
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Kepler’s Third Law: Kepler’s third law states that the square of the period of the orbit of a planet about the Sun is proportional to
the cube of the semi-major axis of the orbit.
The constant of proportionality is
2 2
P 2
planet
P yr
earth
= =1 (5.6.17)
3 3 3
a a AU
planet earth
The net external force on mass m is gravity, and so we substitute the force of gravity for Fnet:
2
mM v
G =m (5.6.19)
r2 r
The fact that m cancels out is another aspect of the oft-noted fact that at a given location all masses fall with the same acceleration.
Here we see that at a given orbital radius r, all masses orbit at the same speed. This was implied by the result of the preceding
worked example. Now, to get at Kepler’s third law, we must get the period P into the equation. By definition, period P is the time
for one complete orbit. Now the average speed v is the circumference divided by the period—that is,
2πr
v = . (5.6.21)
P
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2 3
2
4π a
P = . (5.6.23)
GM
Using subscripts 1 and 2 to denote two different satellites, and taking the ratio of the last equation for satellite 1 to satellite 2 yields
2 3
P r
1 1
= (5.6.24)
2 3
P r
2 2
This is Kepler’s third law. Note that Kepler’s third law is valid only for comparing satellites of the same parent body, because only
then does the mass of the parent body M cancel.
2 3
r
3
for the ratio r
2
. We obtain a relationship that can be used to
P
determine the mass M of a parent body from the orbits of its satellites:
2 3
4π r
M = (5.6.25)
GP2
If r and P are known for a satellite, then the mass M of the parent can be calculated. This principle has been used extensively to find
3
the masses of heavenly bodies that have satellites. Furthermore, the ratio should be a constant for all satellites of the same
r
2
T
3
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Understanding Kepler’s 3 Laws and Orbits: In this video you will be introduced to Kepler’s 3 laws and see how they are
relevant to orbiting objects.
Orbital Maneuvers
An orbital maneuver is the use of propulsion systems to change the orbit of a spacecraft (the rest of the flight is called “coasting”).
learning objectives
Explain purpose of an orbital maneuver
Orbital Maneuvers
In spaceflight, an orbital maneuver is the use of propulsion systems to change the orbit of a spacecraft. The rest of the flight,
especially in a transfer orbit, is called coasting.
Rocket Equation
The Tsiolkovsky rocket equation or ideal rocket equation is an equation useful for considering vehicles that follow the basic
principle of a rocket: a device that can apply acceleration to itself (a thrust) by expelling part of its mass with high speed and
moving due to the conservation of momentum. Specifically, it is a mathematical equation relating the delta-v with the effective
exhaust velocity and both the initial and final mass of a rocket (or other reaction engine).
For any such maneuver (or journey involving a number of such maneuvers):
m0
Δv = ve ln( ), (5.6.26)
m1
where:
m0 is the initial total mass, including propellant;
m1 is the final total mass;
ve is the effective exhaust velocity (ve=Isp • g0 where Isp is the specific impulse expressed as a time period and g0 is the
gravitational constant); and
Δv is delta-v the maximum change of speed of the vehicle (with no external forces acting).
See for an illustration plotting the relationship between final velocity and rocket mass ratios (according to the rocket equation).
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Rocket Equation: Rocket mass ratios versus final velocity calculated from the rocket equation
Delta-v Budget:
The total delta-v for each maneuver estimated for a mission is called a delta-v budget. With a good approximation of the delta-v
budget, designers can estimate the fuel to payload requirements of the spacecraft using the rocket equation.
Voyager Path Using Gravity Assists: The trajectories that enabled NASA’s twin Voyager spacecraft to tour the four gas giant
planets and achieve velocity to escape our solar system
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Transfer Orbits
Orbit insertion is a general term used for a maneuver when it is more than a small correction. It may be used in a maneuver to
change a transfer orbit or an ascent orbit into a stable one, but also to change a stable orbit into a descent (i.e., descent orbit
insertion). Also, the term orbit injection is used, especially for changing a stable orbit into a transfer orbit—e.g., trans-lunar
injection (TLI), trans-Mars injection (TMI) and trans-Earth injection (TEI).
The Hohmann transfer orbit is an elliptical orbit used to transfer between two circular orbits of different altitudes in the same plane.
The orbital maneuver to perform the Hohmann transfer uses two engine impulses that move aspacecraft onto and off the transfer
orbit, as diagramed in. Hohmann transfer orbits are the most efficient with fuel. Other non-Hohmann types of transfer orbits that
are less efficient with fuel exist, but these may be more efficient with other resources (such as time).
Satellites
Natural satellites are celestial objects that orbit a larger body; artificial satellites are manmade objects put in the orbit of the Earth.
learning objectives
Define the concept of a satellite, in the broadest possible terms
Satellites
The word “satellite” has a somewhat ambiguous definition. The broadest possible definition of a satellite is an object that orbits a
larger one due to the force of gravity. Natural satellites, often called moons (see ), are celestial bodies that orbit a larger body call a
primary (often planet, though there are binary asteroids, too). It is technically correct to refer to a planet as a “satellite” of its parent
star, though this is not common.
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Moons of the Solar System: Nineteen natural satellites are large enough to be round, and one, Saturn’s Titan, has a substantial
atmosphere.
Artificial satellites (see ) are man made objects put in orbit about the Earth or another planet in the Solar System. All satellites
follow the laws of orbital mechanics, which can almost always be approximated with Newtonian physics.
Natural Satellites
Formally classified natural satellites, or moons, include 176 planetary satellites orbiting six of the eight planets, and eight orbiting
three of the five IAU-listed dwarf planets. As of January 2012, over 200 minor-planet moons have been discovered. There are 76
objects in the asteroid belt with satellites (five with two satellites each), four Jupiter trojans, 39 near-Earth objects, and 14 Mars-
crossers. There are also 84 known natural satellites of trans-Neptunian objects. Planets around other stars are likely to have natural
satellites as well, although none have yet been observed.
Of the inner planets, Mercury and Venus have no natural satellites; Earth has one large natural satellite, known as the Moon; and
Mars has two tiny natural satellites, Phobos and Deimos. The large gas giants have extensive systems of natural satellites, including
half a dozen comparable in size to Earth’s Moon: the four Galilean moons, Saturn’s Titan, and Neptune’s Triton. Saturn has an
additional six mid-sized natural satellites massive enough to have achieved hydrostatic equilibrium, and Uranus has five. It has
been suggested that some satellites may potentially harbor life, though there is currently no direct evidence.
The Earth–Moon system is unique in that the ratio of the mass of the Moon to the mass of the Earth is much greater than that of any
other natural satellite to planet ratio in the Solar System. Additionally the Moon’s orbit with respect to the Sun is always concave.
The seven largest natural satellites in the Solar System (those bigger than 2,500 km across) are Jupiter’s Galilean moons
(Ganymede, Callisto, Io, and Europa), Saturn’s moon Titan, Earth’s moon, and Neptune’s captured natural satellite Triton.
Artificial Satellites
The first satellite, Sputnik 1, was put into orbit around Earth and was therefore in geocentric orbit. By far this is the most common
type of orbit with about 2,500 artificial satellites orbiting the Earth. Geocentric orbits may be further classified by their altitude,
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inclination and eccentricity.
The commonly used altitude classifications are Low Earth orbit (LEO), Medium Earth orbit (MEO) and High Earth orbit (HEO).
Low Earth orbit is any orbit below 2000 km, and Medium Earth orbit is any orbit higher than that but still below the altitude for
geosynchronous orbit at 35,786 km. High Earth orbit is any orbit higher than the altitude for geosynchronous orbit.
Altitude classifications
Low Earth orbit (LEO): Geocentric orbits ranging in altitude from 0–2000 km (0–1240 miles)
Medium Earth orbit (MEO): Geocentric orbits ranging in altitude from 2,000 km (1,200 mi) to just below geosynchronous orbit
at 35,786 km (22,236 mi). Also known as an intermediate circular orbit.
High Earth orbit (HEO): Geocentric orbits above the altitude of geosynchronous orbit 35,786 km (22,236 mi).
Inclination Classifications
Inclined orbit: An orbit whose inclination in reference to the equatorial plane is not zero degrees.
Polar orbit: An orbit that passes above or nearly above both poles of the planet on each revolution. Therefore it has an
inclination of (or very close to) 90 degrees.
Polar sun synchronous orbit: A nearly polar orbit that passes the equator at the same local time on every pass. Useful for image
taking satellites because shadows will be nearly the same on every pass.
Eccentricity Classifications
Circular orbit: An orbit that has an eccentricity of 0 and whose path traces a circle.
Hohmann transfer orbit: An orbital maneuver that moves a spacecraft from one circular orbit to another using two engine
impulses.
Elliptic orbit: An orbit with an eccentricity greater than 0 and less than 1 whose orbit traces the path of an ellipse.
Geosynchronous transfer orbit: An elliptic orbit where the perigee is at the altitude of a Low Earth orbit (LEO) and the apogee
at the altitude of a geosynchronous orbit.
Geostationary transfer orbit: An elliptic orbit where the perigee is at the altitude of a Low Earth orbit (LEO) and the apogee at
the altitude of a geostationary orbit.
Key Points
An ellipse is a closed plane curve that resembles a stretched out circle (The Sun is at one focus while the other focus has no
physical significance. A circle is a special case of an ellipse where both focal points coincide.
How stretched out an ellipse is from a perfect circle is known as its eccentricity: a parameter that can take any value greater
than or equal to 0 (a circle) and less than 1 (as the eccentricity tends to 1, the ellipse tends to a parabola).
p
Symbolically, an ellipse can be represented in polar coordinates as: r = , where (r, θ) are the polar coordinates (from
1+ϵ cos θ
the focus) for the ellipse, p is the semi-latus rectum, and ϵ is the eccentricity of the ellipse.
p
Perihelion is minimum distance from the Sun a planet achieves in its orbit and is given byr = . Aphelion is the largest
max
1−ϵ
p
distance from the Sun a planet reaches in his orbit and is given by r max = .
1 −ϵ
In a small time the planet sweeps out a small triangle having base line and height. The area of this triangle is given by
1 dθ
dA =
1
2
⋅ r ⋅ rdθ and so the constant areal velocity is: dA
dt
=
2
r
2 dt
The period P satisfies: πab = P ⋅ r θ˙ . One can see that the product of r and must be constant, so that when the planet is
1
2
2 2
further from the Sun it travels at a slower rate and vise versa.
A planet travels fastest at perihelion and slowest at aphelion.
Kepler’s third law can be represented symbolically as P ∝ a , where P is the orbital period of the planet and a is the semi-
2 3
Kepler’s third law can be derived from Newton’s laws of motion and the universal law of gravitation. Set the force of gravity
2
equal to the centripetal force. After substituting an expression for the velocity of the planet, one can obtain: G = which M
r
4πr
2
P
2 3
GM
.
2 3
Using the expression above we can obtain the mass of the parent body from the orbits of its satellites: M = 4π r
2
GP
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The ideal rocket equation related the maximum change in velocity attainable by a rocket ( delta-v or Δv) as a function of the
exhaust velocity (ve) and the ratio between the mass of the rocket with and without the propellant (m0/m1). The equation is
m0
given by Δv = v e ln( ) .
m1
The Oberth effect: where the use of a rocket engine travelling at high speed generates more useful energy than one at low speed.
Thus it is more efficient to apply thrust when the spacecraft is nearest to the planet (periastron).
A gravity assist maneuver is the use of the relative movement and gravity of a planet (or other celestial body) to alter the
velocity of a spacecraft—typically in order to save propellant, time, and expense. This technique was employed by the Voyager
probes (see.
The Hohmann transfer orbit is an elliptical orbit used to transfer between two circular orbits of different altitudes, in the same
plane. The orbital maneuver to perform the Hohmann transfer uses two engine impulses which move aspacecraft onto and off
the transfer orbit. See.
The broadest possible definition of a satellite is an object that orbits a larger one due to the force of gravity.
All satellites follow the laws of orbital mechanics, which can almost always be approximated with Newtonian physics.
Natural satellites are often classified in terms of their size and composition, while artificial satellites are categorized in terms of
their orbital parameters.
Artificial Earth-orbiting satellites have orbits categorized by their altitudes, inclinations, and eccentricities.
Key Terms
rmax − rmin
eccentricity: The coefficient of variation between r min and r
max :ϵ= The further appart the foci are, the stronger
rmax + rmin
the eccentricity.
perihelion: The point in the elliptical orbit of a planet or comet etc. where it is nearest to the Sun. The point farthest from the
Sun is called aphelion.
semi-latus rectum: The latus rectum is a chord perpendicular to the major axis and passing through the focus. The semi-latus
rectum is half the latus rectrum. See distance p in.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
mean motion: An angle of 2π (radians) divided by the orbital period (of a celestial body in an elliptic orbit).
astronomical unit: The mean distance from the Earth to the Sun (the semi-major axis of Earth’s orbit), approximately
149,600,000 kilometres (symbol AU), used to measure distances in the solar system.
sidereal year: The orbital period of the Earth; a measure of the time it takes for the Sun to return to the same position with
respect to the stars of the celestial sphere. A sidereal year is about 20.4 minutes longer than the tropical year due to precession
of the equinoxes.
Hohmann transfer orbit: The Hohmann transfer orbit is an elliptical orbit used to transfer between two circular orbits of
different altitudes, in the same plane. The orbital maneuver to perform the Hohmann transfer uses two engine impulses, one to
move a spacecraft onto the transfer orbit and a second to move off it.
delta-v: The maximum change in the scalar speed of a rocket if the rocket were operated in a vacuum away from external
forces (i.e., if no other external forces act).
natural satellite: A natural satellite, moon, or secondary planet is a celestial body that orbits a planet or smaller body, which is
called its primary.
artificial satellite: In the context of spaceflight, a satellite is an object which has been placed into orbit by human endeavour.
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5.7: Gravitational Potential Energy
Defining Graviational Potential Energy
Gravitational energy is the potential energy associated with gravitational force, such as elevating objects against the Earth’s gravity.
learning objectives
Express gravitational potential energy for two masses
Gravitational energy is the potential energy associated with gravitational force, as work is required to elevate objects against
Earth’s gravity. The potential energy due to elevated positions is called gravitational potential energy, and is evidenced by water in
an elevated reservoir or kept behind a dam. If an object falls from one point to another point inside a gravitational field, the force of
gravity will do positive work on the object, and the gravitational potential energy will decrease by the same amount.
Consider a book placed on top of a table. As the book is raised from the floor to the table, some external force works against the
gravitational force. If the book falls back to the floor, the “falling” energy the book receives is provided by the gravitational force.
Thus, if the book falls off the table, this potential energy goes to accelerate the mass of the book and is converted into kinetic
energy. When the book hits the floor, this kinetic energy is converted into heat and sound by the impact.
The factors that affect an object’s gravitational potential energy are its height relative to some reference point, its mass, and the
strength of the gravitational field it is in. Thus, a book lying on a table has less gravitational potential energy than the same book on
top of a taller cupboard, and less gravitational potential energy than a heavier book lying on the same table. An object at a certain
height above the Moon’s surface has less gravitational potential energy than at the same height above the Earth’s surface because
the Moon’s gravity is weaker. Note that “height” in the common sense of the term cannot be used for gravitational potential energy
calculations when gravity is not assumed to be a constant. The following sections provide more detail.
Local Approximation
The strength of a gravitational field varies with location. However, when the change of distance is small in relation to the distances
from the center of the source of the gravitational field, this variation in field strength is negligible and we can assume that the force
of gravity on a particular object is constant. Near the surface of the Earth, for example, we assume that the acceleration due to
gravity is a constant g = 9.8m/s (“standard gravity”). In this case, a simple expression for gravitational potential energy can be
2
derived using the W = Fd equation for work. The upward force required while moving at a constant velocity is equal to the
weight, mg, of an object, so the work done in lifting it through a height h is the product mgh. Thus, when accounting only for
mass, gravity, and altitude, the equation is:
U = mgh (5.7.1)
where U is the potential energy of the object relative to its being on the Earth’s surface, m is the mass of the object, g is the
acceleration due to gravity, and h is the altitude of the object. If m is expressed in kilograms, g in m/s and h in meters then U
2
will be calculated in joules. In most situations, the change in potential energy is the relevant quantity:
ΔU = mgΔh (5.7.2)
General Formula
However, over large variations in distance, the approximation that gg is constant is no longer valid, and we have to use calculus and
the general mathematical definition of work to determine gravitational potential energy. For the computation of the potential
energy, we can integrate the gravitational force, whose magnitude is given by Newton’s law of gravitation, with respect to the
distance r between the two bodies. Using that definition, the gravitational potential energy of a system of masses m and M at a
1 2
where K is the constant of integration. Choosing the convention that K = 0 makes calculations simpler, albeit at the cost of making
U negative. Note that in this case the potential energy becomes zero when r is infinite, and approaches negative infinity as r goes
to zero.
5.7.1 https://phys.libretexts.org/@go/page/15449
Trebuchet: A trebuchet uses the gravitational potential energy of the counterweight to throw projectiles over long distances.
Key Points
If an object falls from one point to another point inside a gravitational field, the force of gravity will do positive work on the
object, and the gravitational potential energy will decrease by the same amount.
Near the surface of the Earth, the work done in lifting an object through a height h is the product mgh, so U = mgh .
The gravitational potential energy, U , of a system of masses m and M at a distance r using gravitational constant G is
1 2
m1 M2
U = −G +K .
r
Key Terms
Newton’s law of gravitation: This law states that every point mass in the universe attracts every other point mass with a force
that is directly proportional to the product of their masses and inversely proportional to the square of the distance between them.
potential energy: The energy an object has because of its position (in a gravitational or electric field) or its condition (as a
stretched or compressed spring, as a chemical reactant, or by having rest mass)
gravity: Resultant force on Earth’s surface, of the attraction by the Earth’s masses, and the centrifugal pseudo-force caused by
the Earth’s rotation.
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5.8: Energy Conservation
learning objectives
Calculate the escape speed of an object given its kinetic energy and the gravitational potential energy
Escape speed is the required starting speed required by an object to go from a starting point in a gravitational potential field to an
ending point that is infinitely far away. It is assumed that the velocity of the object at the ending point will be zero.
Isaac Newton’s Analysis of Escape Speed: In this figure, Objects A and B don’t have the required escape speed and so they fall
back to Earth after launch. Objects C and D don’t either, they achieve a circular and an elliptical orbit respectively. Object E is
launched with sufficient escape velocity and escapes the Earth.
Imagine a situation in which a spaceship that does not have a propulsion system is launched straight away from a planet. (It is moot
to discuss escape speed for objects with propulsion systems.) Let us assume that the only significant force that is acting on the
spaceship is the force of gravity from the planet. The escape speed of the spaceship can calculated through a simple analysis of
conservation of energy. The gravitational potential energy of the spaceship is:
GMm
U =− (5.8.1)
r
Where G is the universal gravitational constant (G = 6.67 ⋅ 10 m kg s ), M is the mass of the planet,
−11 3 −1 −2
m is the mass of the
spaceship, and r is the distance of the spaceship from the planet’s center of gravity.
At the ending point of the spaceship, r goes to infinity. As r goes to infinity, the value of the gravitational potential energy
expression goes to 0.
The kinetic energy of the spaceship can be found from:
1
2
mv (5.8.2)
2
Where mm is the mass of the spaceship and v is the velocity of the spaceship.
At the starting point of the spaceship, the velocity must have a magnitude equal to the escape speed (s ). The velocity of the
e
spaceship is 0 at its ending point, and so consequently its kinetic energy is 0 in the end as well.
Summarizing the kinetic energy (K) and potential energy (U ) of the spaceship at it’s initial (i) and final (f ) states:
1 −GMm
2
(K + U)i = mse + (5.8.3)
2 r
(K + U)f = 0 + 0 (5.8.4)
Due to conservation of energy, the initial energy must equal the final energy and so we can solve for s : e
−−−−−
2GM
se =√ (5.8.5)
r
Interestingly, if the spaceship were to fall to the planet from a point infinitely far away it would obtain a final speed of se at the
planet.
5.8.1 https://phys.libretexts.org/@go/page/15450
It should be noted that if an object is launched from a rotating body, such as the Earth, the speed at which the body rotates will
affect the required velocity that an object must have relative to the surface of the body. If a rocket is launched tangentially from the
Earth’s equator in the same direction that the Earth is turning, it will require a lower velocity relative to the Earth than if it were
launched in the opposite direction to meet escape speed requirements.
Additionally, it is a misconception that powered vehicles (such as rockets) require escape speed to leave orbit and travel through
outer-space. If the vehicle has a propulsion system to provide it with energy once it has left the surface of the planet, it is not
necessary to initially meet escape speed requirements.
Key Points
It is assumed that the velocity of the object at the ending point will be zero.
−−−−−
2GM
The requisite escape speed (se)(se) of an object to escape a spherically symmetric body is given by: s e =√ , where G is
r
the universal gravitational constant, M is the mass of the body, and r is the distance of the object from the body’s center of
gravity.
Escape speed is the required speed that an object has to have to go from a starting point in a gravitational potential field to an
ending point that is infinitely far away.
The speed at which a body rotates will affect the required velocity that an object must have relative to the surface of the body.
Objects that have propulsion systems do not need to reach escape velocity.
Key Terms
propulsion: Force causing movement.
potential energy: The energy an object has because of its position (in a gravitational or electric field) or its condition (as a
stretched or compressed spring, as a chemical reactant, or by having rest mass)
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
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5.9: Angular vs. Linear Quantities
learning objectives
Describe properties characteristic to angular velocity and angular momentum
Linear motion is motion in a straight line. This type of motion has several familiar vector quantities associated with it, including
linear velocity and momentum. These vector quantities each have a magnitude (a scalar, or number) and direction associated with
them. Similarly, circular motion is motion in a circle. It has the same set of vector quantities associated with it, including angular
velocity and angular momentum.
Angular velocity diagram: A vector diagram illustrating circular motion. The blue vector connects the origin (center) of the
motion to the position of the particle. The red vector is the angular velocity vector, pointing perpendicular to the plane of motion
and with magnitude equal to the instantaneous velocity. File:Angular velocity.svg - Wikipedia, the free encyclopedia. Provided by:
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Imagine a particle moving in a circle around a point at a constant speed. We will call that point the origin. At any instant in time,
the particle is moving in a particular straight-line direction with that speed. In the next instant, the particle has the same speed, but
the direction of its velocity has changed.
We recall from our study of linear velocity that a change in the direction of the velocity vector, is a change in velocity and a change
in velocity is acceleration. However, we can define an angular momentum vector which is constant throughout this motion. The
angular velocity has a direction perpendicular to the plane of circular motion, just like a bike axle points perpendicularly to the
rotating wheel. This direction never changes as the object moves in its circle. The magnitude of the angular momentum is equal to
the rate at which the angle of the particle advances:
dϕ
ω = (5.9.1)
dt
Note that there are two vectors that are perpendicular to any plane. For example, imagine a vector pointing into your table and the
opposite one pointing out of it. To remove this ambiguity, the convention in physics is to use the right hand rule: curl the fingers of
your right hand in the direction of the circular motion, and your thumb will point toward the direction of the angular velocity and
momentum vectors.
Right hand rule: When determining the direction of an angular vector, use the right hand rule: curl the fingers of your right hand
in the direction of the circular motion and your thumb points in the vector direction. File:Right-hand grip rule.svg - Wikipedia, the
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The units of angular velocity are radians per second. Radian describes the plane angle subtended by a circular arc as the length of
the arc divided by the radius of the arc. One radian is the angle subtended at the center of a circle by an arc that is equal in length to
the radius of the circle. More generally, the magnitude in radians of such a subtended angle is equal to the ratio of the arc length to
the radius of the circle; that is, θ = , where θ is the subtended angle in radians, s is arc length, and r is radius.
s
Thus, while the object moves in a circle at constant speed, it undergoes constant linear acceleration to keep it moving in a circle.
However, it’s angular velocity is constant since it continually sweeps out a constant arc length per unit time. Constant angular
velocity in a circle is known as uniform circular motion.
Just as there is an angular version of velocity, so too is there an angular version of acceleration. When the object is going around a
circle but its speed is changing, the object is undergoing angular acceleration. Just like with linear acceleration, angular acceleration
is a change in the angular velocity vector. This change could be a change in the speed of the object or in the direction. Angular
velocity can be clockwise or counterclockwise.
Key Points
The direction of angular quantity vectors points perpendicular to the plane of the motion. You can determine this direction using
the right hand rule.
The direction of linear quantities such as velocity and momentum change as an object moves in a circle. We can instead define
angular versions of these quantities which are constant throughout the circular motion.
The units of angular quantities are per radian, a measurement of angle, rather than per linear distance (e.g. meter).
Key Terms
vector: A directed quantity, one with both magnitude and direction; the between two points.
angular momentum: A vector quantity describing an object in circular motion; its magnitude is equal to the momentum of the
particle, and the direction is perpendicular to the plane of its circular motion.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
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CHAPTER OVERVIEW
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1
6.1: Introduction
Introduction to Work and Energy
Work is the energy associated with the action of a force.
learning objectives
Describe relationship between work, energy, and force
The work done on a system by a constant force is the product of the component of the force in the direction of motion times the
distance through which the force acts. For one-way motion in one dimension, this is expressed in equation form as W = Fd cos θ ,
where W is work, F is the magnitude of the force on the system, d is the magnitude of the displacement of the system, and θ is the
angle between the force vector F and the displacement vector d.
because there is no motion. No energy is transferred to or from the briefcase. (C) The person moving the briefcase horizontally at a
constant speed does no work on it, and transfers no energy to it. (D) Work is done on the briefcase by carrying it up stairs at
constant speed, because there is necessarily a component of force F in the direction of the motion. Energy is transferred to the
briefcase and could, in turn, be used to do work. (E) When the briefcase is lowered, energy is transferred out of the briefcase and
6.1.1 https://phys.libretexts.org/@go/page/14466
into an electric generator. Here the work done on the briefcase by the generator is negative, removing energy from the briefcase,
because F and d are in opposite directions.
Examples of Work: This is how work in progress and energy co-exist and operate. Work is the energy associated with the action of
a force.
In physics, a force is said to do work when it acts on a body so that there is a displacement of the point of application, however
small, in the direction of the force. Thus a force does work when there is movement under the action of the force. The work done
by a constant force of magnitude F on a point that moves a distance d in the direction of the force is the product:
W = Fd (6.1.1)
For example, if a force of 10 newton (F = 10 N) acts along point that travels two meters (d = 2 m), then it does the work W = (10
N)(2 m) = 20 N m = 20 J. This is approximately the work done lifting a 1 kg weight from the ground to over a person’s head
against the force of gravity. Notice that the work is doubled either by lifting twice the weight in the same distance or by lifting the
same weight twice the distance.
Work is closely related to energy. The conservation of energy states that the change in total internal energy of a system equals the
added heat minus the work performed by the system (see the first law of thermodynamics, and ):
Baseball Pitcher: A baseball pitcher does work on a baseball by throwing the ball at some force, F, over some distance d, which
for the average baseball field, is about 60 feet.
6.1.2 https://phys.libretexts.org/@go/page/14466
δE = δQ − δW (6.1.2)
Also, from Newton’s second law for rigid bodies, it can be shown that work on an object is equal to the change in kinetic energy of
that object:
W = ΔKE (6.1.3)
The work of forces generated by a potential function is known as potential energy and the forces are said to be conservative.
Therefore work on an object moving in a conservative force field is equal to minus the change of potential energy of the object:
W = −ΔPE (6.1.4)
This shows that work is the energy associated with the action of a force, and so has the physical dimensions and units of energy.
Key Points
Work is the transfer of energy by a force acting on an object as it is displaced.
The work done by a force is zero if the displacement is either zero or perpendicular to the force.
The work done is positive if the force and displacement have the same direction, and the work done is negative if they have
opposite direction.
Key Terms
energy: A quantity that denotes the ability to do work and is measured in a unit dimensioned in mass × distance²/time² (ML²/T²)
or the equivalent.
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6.2: Work Done by a Constant Force
Force in the Direction of Displacement
The work done by a constant force is proportional to the force applied times the displacement of the object.
learning objectives
Contrast displacement and distance in constant force situations
Units
One way to validate if an expression is correct is to perform dimensional analysis. We have claimed that work is the change in
kinetic energy of an object and that it is also equal to the force times the distance. The units of these two should agree. Kinetic
energy – and all forms of energy – have units of joules (J). Likewise, force has units of newtons (N) and distance has units of
meters (m). If the two statements are equivalent they should be equivalent to one another.
2
m m
N ⋅ m = kg ⋅ m = kg =J (6.2.1)
2 2
s s
Example Problems
Here are a few example problems:
(1.a) Consider a constant force of two newtons (F = 2 N) acting on a box of mass three kilograms (M = 3 kg). Calculate the work
done on the box if the box is displaced 5 meters.
(1.b) Since the box is displaced 5 meters and the force is 2 N, we multiply the two quantities together. The object’s mass will
dictate how fast it is accelerating under the force, and thus the time it takes to move the object from point a to point b. Regardless
of how long it takes, the object will have the same displacement and thus the same work done on it.
(2.a) Consider the same box (M = 3 kg) being pushed by a constant force of four newtons (F = 4 N). It begins at rest and is pushed
for five meters (d = 5m). Assuming a frictionless surface, calculate the velocity of the box at 5 meters.
(2.b) We now understand that the work is proportional to the change in kinetic energy, from this we can calculate the final velocity.
What do we know so far? We know that the block begins at rest, so the initial kinetic energy must be zero. From this we
algebraically isolate and solve for the final velocity.
1
2
Fd = ΔKE = KEf − 0 = mv (6.2.2)
f
2
−−−− −−−−− − −−
Fd 4N ⋅ 5m −−
vf = √2 = √2 = √10m/s (6.2.3)
m 2kg
We see that the final velocity of the block is approximately 3.15 m/s.
6.2.1 https://phys.libretexts.org/@go/page/14467
Force at an Angle to Displacement
A force does not have to, and rarely does, act on an object parallel to the direction of motion.
learning objectives
Infer how to adjust one-dimensional motion for our three-dimensional world
The Fundamentals
Up until now, we have assumed that any force acting on an object has been parallel to the direction of motion. We have considered
our motion to be one dimensional, only acting along the x or y axis. To best examine and understand how nature operators in our
three-dimensional world, we will first discuss work in two dimensions in order to build our intuition.
A force does not have to, and rarely does, act on an object parallel to the direction of motion. In the past, we derived that W = Fd ;
such that the work done on an object is the force acting on the object multiplied by the displacement. But this is not the whole
story. This expression contains an assumed cosine term, which we do not consider for forces parallel to the direction of motion.
“Why would we do such a thing? ” you may ask. We do this because the two are equivalent. If the angle of the force along the
direction of motion is zero, such that the force is parallel to the direction of motion, then the cosine term equals one and does not
change the expression. As we increase the force’s angle with respect to the direction of motion, less and less work is done along the
direction that we are considering; and more and more work is being done in another, perpendicular, direction of motion. This
process continues until we are perpendicular to our original direction of motion, such that the angle is 90, and the cosine term
would equal zero; resulting in zero work being done along our original direction. Instead, we are doing work in another direction!
Angle: Recall that both the force and direction of motion are vectors. When the angle is 90 degrees, the cosine term goes to zero.
When along the same direction, they equal one.
Let’s show this explicitly and then look at this phenomena in terms of a box moving along the x and y directions.
We have discussed that work is the integral of the force and the dot product respect to x. But in fact, dot product of force and a very
small distance is equal to the two terms times cosine of the angle between the two. F × dx = Fd cos(θ) . Explicitly,
t1 t1
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In the second scenario, the box is being pushed at an angle of 45 degrees to the x-direction; and thus also a 45 degree angle to the
y-direction. When evaluated, the cosine of 45 degrees is equal to , or approximately 0.71. This means is that 71% of the force is
1
√2
contributing to the work along the x-direction. The other 29% is acting along the y-direction.
In the third scenario, we know that the force is acting at a 60 degree angle to the x-direction; and thus also a 30 degree angle to the
y-direction. When evaluated, cosine of 60 degrees is equal to 1/2. This means that the force is equally acting in the x and y-
direction! The work done is linear with respect to both x and y.
In the last scenario, the box is being pushed at an angle perpendicular to the x direction. In other words, we are pushing the box in
the y-direction! Thus, the box’s position will be unchanged and experience no displacement along the x-axis. The work done in the
x direction will be zero.
Key Points
Understanding work is quintessential to understanding systems in terms of their energy, which is necessary for higher level
physics.
Work is equivalent to the change in kinetic energy of a system.
Distance is not the same as displacement. If a box is moved 3 meters forward and then 4 meters to the left, the total
displacement is 5 meters, not 7 meters.
Work done on an object along a given direction of motion is equal to the force times the displacement times the cosine of the
angle.
No work is done along a direction of motion if the force is perpendicular.
When considering force parallel to the direction of motion, we omit the cosine term because it equals 1 which does not change
the expression.
Key Terms
work: A measure of energy expended in moving an object; most commonly, force times displacement. No work is done if the
object does not move.
dot product: A scalar product.
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6.3: Work Done by a Variable Force
Work Done by a Variable Force
Integration is used to calculate the work done by a variable force.
learning objectives
Describe approaches used to calculate work done by a variable force
W = F ⋅ Δx (6.3.1)
In the case of a variable force, integration is necessary to calculate the work done. For example, let’s consider work done by a
spring. According to the Hooke’s law the restoring force (or spring force) of a perfectly elastic spring is proportional to its
extension (or compression), but opposite to the direction of extension (or compression). So the spring force acting upon an object
attached to a horizontal spring is given by:
Fs = −kx (6.3.2)
that is proportional to its displacement (extension or compression) in the x direction from the spring’s equilibrium position, but its
direction is opposite to the x direction. For a variable force, one must add all the infinitesimally small contributions to the work
done during infinitesimally small time intervals dt (or equivalently, in infinitely small length intervals dx=vxdt). In other words, an
integral must be evaluated:
t t x
1
2
Ws = ∫ Fs ⋅ vdt = ∫ −kxvx dt = ∫ −kxdx = − kΔx (6.3.3)
0 0 xo
2
This is the work done by a spring exerting a variable force on a mass moving from position xo to x (from time 0 to time t). The
work done is positive if the applied force is in the same direction as the direction of motion; so the work done by the object on
spring from time 0 to time t, is:
t t
1
2
Wa = ∫ Fa ⋅ vdt = ∫ −Fs ⋅ vdt = kΔx (6.3.4)
0 0
2
in this relation F is the force acted upon spring by the object. F and F are in fact action- reaction pairs; and W is equal to the
a a s a
Another example is the work done by gravity (a constant force) on a free-falling object (we assign the y-axis to vertical motion, in
this case):
t2 t2 y2
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Notice that the result is the same as we would have obtained by simply evaluating the product of force and distance.
Key Points
The work done by a constant force of magnitude F on a point that moves a displacement d in the direction of the force is the
product: W = Fd .
Integration approach can be used both to calculate work done by a variable force and work done by a constant force.
The SI unit of work is the joule; non- SI units of work include the erg, the foot-pound, the foot-poundal, the kilowatt hour, the
litre-atmosphere, and the horsepower-hour.
Key Terms
work: A measure of energy expended in moving an object; most commonly, force times displacement. No work is done if the
object does not move.
force: A physical quantity that denotes ability to push, pull, twist or accelerate a body, which is measured in a unit dimensioned
in mass × distance/time² (ML/T²): SI: newton (N); CGS: dyne (dyn)
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6.4: Work-Energy Theorem
Kinetic Energy and Work-Energy Theorem
The work-energy theorem states that the work done by all forces acting on a particle equals the change in the particle’s kinetic
energy.
learning objectives
Outline the derivation of the work-energy theorem
Kinetic Energy: A force does work on the block. The kinetic energy of the block increases as a result by the amount of work. This
relationship is generalized in the work-energy theorem.
The work W done by the net force on a particle equals the change in the particle’s kinetic energy KE:
1 1
2 2
W = ΔKE = mv − mv (6.4.1)
f i
2 2
where vi and vf are the speeds of the particle before and after the application of force, and m is the particle’s mass.
Derivation
For the sake of simplicity, we will consider the case in which the resultant force F is constant in both magnitude and direction and
is parallel to the velocity of the particle. The particle is moving with constant acceleration a along a straight line. The relationship
between the net force and the acceleration is given by the equation F = ma (Newton’s second law), and the particle’s displacement
d, can be determined from the equation:
2 2
v =v + 2ad (6.4.2)
f i
obtaining,
2 2
v −v
f i
d= (6.4.3)
2a
The work of the net force is calculated as the product of its magnitude (F=ma) and the particle’s displacement. Substituting the
above equations yields:
2 2
v −v 1 1
f i 2 2
W = Fd = ma = mv − mv = KEf − KEi = ΔKE (6.4.4)
f i
2a 2 2
Key Points
The work W done by the net force on a particle equals the change in the particle’s kinetic energy KE:
W = ΔKE =
1
2
mv −
2
f
1
2
mv
2
i
.
The work-energy theorem can be derived from Newton’s second law.
Work transfers energy from one place to another or one form to another. In more general systems than the particle system
mentioned here, work can change the potential energy of a mechanical device, the heat energy in a thermal system, or the
electrical energy in an electrical device.
6.4.1 https://phys.libretexts.org/@go/page/14469
Key Terms
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
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6.5: Potential Energy and Conservation of Energy
Conservative and Nonconservative Forces
Conservative force—a force with the property that the work done in moving a particle between two points is independent of the
path it takes.
learning objectives
Describe properties of conservative and nonconservative forces
A conservative force is a force with the property that the work done in moving a particle between two points is independent of the
path taken. Equivalently, if a particle travels in a closed loop, the net work done (the sum of the force acting along the path
multiplied by the distance travelled) by a conservative force is zero.
A conservative force is dependent only on the position of the object. If a force is conservative, it is possible to assign a numerical
value for the potential at any point. When an object moves from one location to another, the force changes the potential energy of
the object by an amount that does not depend on the path taken. Gravity and spring forces are examples of conservative forces.
If a force is not conservative, then defining a scalar potential is not possible, because taking different paths would lead to
conflicting potential differences between the start and end points. Nonconservative forces transfer energy from the object in motion
(just like conservative force), but they do not transfer this energy back to the potential energy of the system to regain it during
reverse motion. Instead, they transfer the energy from the system in an energy form which can not be used by the force to transfer it
back to the object in motion. Friction is one such nonconservative force.
Motion Along Different Paths: Motion along different paths. For a conservative force, work done via different path is the same.
W = WAB1 + WBA2 = 0. (6.5.1)
Let us now change the path for motion from A to B by another path, shown as path 3. Again, the total work by the conservative
force for the round trip via new route is zero: W = W +W
AB3 = 0. .
BA2
Comparing two equations, W AB1=W . This is true for an arbitrary path. Therefore, work done for motion from A to B by
AB3
learning objectives
Relate the potential energy and the work
Potential energy is often associated with restoring forces such as a spring or the force of gravity. The action of stretching the spring
or lifting the mass of an object is performed by an external force that works against the force field of the potential. This work is
stored in the force field as potential energy. If the external force is removed the force field acts on the body to perform the work as
it moves the body back to its initial position, reducing the stretch of the spring or causing the body to fall. The more formal
6.5.1 https://phys.libretexts.org/@go/page/14470
definition is that potential energy is the energy difference between the energy of an object in a given position and its energy at a
reference position.
Potential Energy in a Bow and Arrow: In the case of a bow and arrow, the energy is converted from the potential energy in the
archer’s arm to the potential energy in the bent limbs of the bow when the string is drawn back. When the string is released, the
potential energy in the bow limbs is transferred back through the string to become kinetic energy in the arrow as it takes flight.
If the work for an applied force is independent of the path, then the work done by the force is evaluated at the start and end of the
trajectory of the point of application. This means that there is a function U(x), called a “potential,” that can be evaluated at the two
points x(t = t1) and x(t2) to obtain the work over any trajectory between these two points. It is tradition to define this function with a
negative sign so that positive work is represented as a reduction in the potential:
x( t2 )
W =∫ F ⋅ dx = ∫ F ⋅ dx (6.5.2)
C x( t1 )
Gravity
Gravitational energy is the potential energy associated with gravitational force, as work is required to move objects against gravity.
learning objectives
Generate an equation that can be used to express the gravitational potential energy near the earth
Gravitational energy is the potential energy associated with gravitational force (a conservative force), as work is required to elevate
objects against Earth’s gravity. The potential energy due to elevated positions is called gravitational potential energy, evidenced, for
example, by water held in an elevated reservoir or behind a dam (as an example, shows Hoover Dam). If an object falls from one
point to another point inside a gravitational field, the force of gravity will do positive work on the object, and the gravitational
potential energy will decrease by the same amount.
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Hoover Dam: Hoover dam uses the stored gravitational potential energy to generate electricity.
where PE = potential energy measured in joules (J), m = mass of the object (measured in kg), and h = perpendicular height from the
reference point (measured in m); g = gravitational acceleration (9.8m/s2). Near the surface of the Earth, g can be considered
constant.
General Formula
However, over large variations in distance, the approximation that g is constant is no longer valid. Instead, we must use calculus
and the general mathematical definition of work to determine gravitational potential energy. For the computation of the potential
energy we can integrate the gravitational force, whose magnitude is given by Newton’s law of gravitation (with respect to the
distance r between the two bodies). Using that definition, the gravitational potential energy of a system of masses m and M at a
distance r using gravitational constant G is:
mM ′
mM
U(r) = ∫ (G )dr = −G + K, (6.5.5)
′2
r r r
where K is the constant of integration. Choosing the convention that K=0 makes calculations simpler, albeit at the cost of making U
negative. For this choice, the potential at infinity is defined as 0.
Springs
When a spring is stretched/compressed from its equilibrium position by x, its potential energy is give as U = 1
2
2
kx .
learning objectives
Explain how potential energy is stored in springs
Spring force is conservative force, given by the Hooke’s law: F = −kx , where k is spring constant, measured experimentally for a
particular spring and x is the displacement. We would like to obtain an expression for the work done to the spring. From the
conservation of mechanical energy (Check our Atom on “Conservation of Mechanical Energy), the work should be equal to the
potential energy stored in spring. The displacement x is usually measured from the position of “neutral length” or “relaxed length”
6.5.3 https://phys.libretexts.org/@go/page/14470
– the length of spring corresponding to situation when spring is neither stretched nor compressed. We shall identify this position as
the origin of coordinate reference (x=0).
Hooke’s Law: Plot of applied force F vs. elongation X for a helical spring according to Hooke’s law (solid line) and what the
actual plot might look like (dashed line). Red is used extension, blue for compression. At bottom, schematic pictures of spring
states corresponding to some points of the plot; the middle one is in the relaxed state (no force applied).
Let x = 0 and x = x (> 0) be the initial and final positions of the block attached to the string. As the block slowly moves, we do
f
xf
work W on the spring: W = ∫ (kx)dx = kx . When we stretch the spring. We have to apply force in the same direction as the
0
1
2
2
f
displacement. (Technically, work is given as the inner product of the two vectors: force and displacement. W = F ⋅ Δx ).
Therefore, the overall sign in the integral is +, not -.
If the block is gently released from the stretched position (\mathrm{x = x_f}\), the stored potential energy in the spring will start to
be converted to the kinetic energy of the block, and vice versa. Neglecting frictional forces, Mechanical energy conservation
demands that, at any point during its motion,
1 1
2 2
Total Energy = mv + kx (6.5.6)
2 2
1
2
= kx = constant. (6.5.7)
f
2
From the energy conservation, we can estimate that, by the time the block reaches x=0 position, its speed will be
−
−
v(x = 0) = √
m
k
xf . The block will keep oscillating between x = -xf and xf.
learning objectives
Formulate the principle of the conservation of the mechanical energy
Conservation of mechanical energy states that the mechanical energy of an isolated system remains constant in time, as long as the
system is free of all frictional forces. In any real situation, frictional forces and other non-conservative forces are always present,
but in many cases their effects on the system are so small that the principle of conservation of mechanical energy can be used as a
fair approximation. An example of a such a system is shown in. Though energy cannot be created nor destroyed in an isolated
system, it can be internally converted to any other form of energy.
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A Mechanical System: An example of a mechanical system: A satellite is orbiting the Earth only influenced by the conservative
gravitational force and the mechanical energy is therefore conserved. This acceleration is represented by a green acceleration vector
and the velocity is represented by a red velocity vector.
Derivation
Let us consider what form the work -energy theorem takes when only conservative forces are involved (leading us to the
conservation of energy principle). The work-energy theorem states that the net work done by all forces acting on a system equals its
change in kinetic energy (KE). In equation form, this is:
1 1
2 2
Wnet = mv − mv = ΔKE. (6.5.8)
0
2 2
If only conservative forces act, then W net = Wc , where Wc is the total work done by all conservative forces. Thus, W c = ΔKE .
Now, if the conservative force, such as the gravitational force or a spring force, does work, the system loses potential energy (PE).
That is, W = −PE . Therefore,
c
This equation means that the total kinetic and potential energy is constant for any process involving only conservative forces. That
is,
KE + PE = const or KEi + PEi = KEf + PEf , (6.5.10)
where i and f denote initial and final values. This equation is a form of the work-energy theorem for conservative forces; it is
known as the conservation of mechanical energy principle.
Remember that the law applies to the extent that all the forces are conservative, so that friction is negligible. The total kinetic plus
potential energy of a system is defined to be its mechanical energy (KE + PE ). In a system that experiences only conservative
forces, there is a potential energy associated with each force, and the energy only changes form between KE and various types of
PE (with the total energy remaining constant).
6.5.5 https://phys.libretexts.org/@go/page/14470
Conservation of Mechanical Energy Worked Exampl…
Exampl…
learning objectives
Identify steps necessary to solve a conservation of energy problem
Problem-solving Strategy
You should follow a series of steps whenever you are problem solving:
Step One
Determine the system of interest and identify what information is given and what quantity is to be calculated. For example, let’s
assume you have the problem with car on a roller coaster. You know that the cars of a roller coaster reach their maximum kinetic
energy (KEKE) when at the bottom of their path. When they start rising, the kinetic energy begins to be converted to gravitational
potential energy (PEgPEg). The sum of kinetic and potential energy in the system should remain constant, if losses to friction are
ignored.
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Determining Energy: The cars of a roller coaster reach their maximum kinetic energy when at the bottom of their path. When they
start rising, the kinetic energy begins to be converted to gravitational potential energy. The sum of kinetic and potential energy in
the system remains constant, ignoring losses to friction.
Step Two
Examine all the forces involved and determine whether you know or are given the potential energy from the work done by the
forces. Then use step three or step four.
Step Three
If you know the potential energies (PE) for the forces that enter into the problem, then forces are all conservative, and you can
apply conservation of mechanical energy simply in terms of potential and kinetic energy. The equation expressing conservation of
energy is:
Step Four
If you know the potential energy for only some of the forces, then the conservation of energy law in its most general form must be
used:
where O stand for all other energies, and W stands for work done by non-conservative forces. In most problems, one or more of
E nc
the terms is zero, simplifying its solution. Do not calculate W , the work done by conservative forces; it is already incorporated in
c
the PE terms.
Step Five
You have already identified the types of work and energy involved (in step two). Before solving for the unknown, eliminate terms
wherever possible to simplify the algebra. For example, choose height h = 0 at either the initial or final point—this will allow to
set PEg at zero. Then solve for the unknown in the customary manner.
Step Six
Check the answer to see if it is reasonable. Once you have solved a problem, reexamine the forms of work and energy to see if you
have set up the conservation of energy equation correctly. For example, work done against friction should be negative, potential
energy at the bottom of a hill should be less than that at the top, and so on.
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Joy physics probs 05- Energy conservation
Energy conservation: Part of a series of videos on physics problem-solving. The problems are taken from “The Joy of Physics. ”
This one deals with energy conservation. The viewer is urged to pause the video at the problem statement and work the problem
before watching the rest of the video.
learning objectives
Express the energy conservation relationship that can be applied to solve problems with dissipative forces
INTRODUCTION
We have seen a problem-solving strategy with the conservation of energy in the previous section. Here we will adopt the strategy
for problems with dissipative forces. Since the work done by nonconservative (or dissipative) forces will irreversibly alter the
energy of the system, the total mechanical energy (KE + PE) changes by exactly the amount of work done by nonconservative
forces (Wc). Therefore, we obtain KE + PE + W = KE + PE , where KE and PE represent kinetic and potential energies
i i nc f f
respectively. Therefore, using the new energy conservation relationship, we can apply the same problem-solving strategy as with
the case of conservative forces.
EXAMPLE
Consider the situation shown in, where a baseball player slides to a stop on level ground. Using energy considerations, calculate
the distance the 65.0-kg baseball player slides, given that his initial speed is 6.00 m/s and the force of friction against him is a
6.5.8 https://phys.libretexts.org/@go/page/14470
constant 450 N.
Fig 1: The baseball player slides to a stop in a distance d. In the process, friction removes the player’s kinetic energy by doing an
amount of work fd equal to the initial kinetic energy.
Strategy: Friction stops the player by converting his kinetic energy into other forms, including thermal energy. In terms of the
work-energy theorem, the work done by friction (f), which is negative, is added to the initial kinetic energy to reduce it to zero. The
work done by friction is negative, because f is in the opposite direction of the motion (that is, θ = 180º , and so cos θ = −1 ). Thus
= −f d . The equation simplifies to mv − f d = 0 .
1 2
W nc
2 i
Solution: Solving the previous equation for d and substituting known values yields, we get d = 2.60 m. The most important point of
this example is that the amount of nonconservative work equals the change in mechanical energy.
Key Points
If a particle travels in a closed loop, the net work done (the sum of the force acting along the path multiplied by the distance
travelled) by a conservative force is zero.
Conservative force is dependent only on the position of the object. If a force is conservative, it is possible to assign a numerical
value for the potential at any point.
Nonconservative force transfer the energy from the system in an energy form which can not be used by the force to transfer
back to the object in motion.
If the work for an applied force is independent of the path, then the work done by the force is evaluated at the start and end of
the trajectory of the point of application. This means that there is a function U(x), called a ” potential “.
It is tradition to define the potential function with a negative sign so that positive work is represented as a reduction in the
potential.
Every conservative force gives rise to potential energy. Examples are elastic potential energy, gravitational potential energy, and
electric potential energy.
Gravitational potential energy near the earth can be expressed with respect to the height from the surface of the Earth as PE =
mgh. g = gravitational acceleration (9.8m/s2). Near the surface of the Earth, g can be considered constant.
Over large variations in distance, the approximation that g is constant is no longer valid and a general formula should be used
for the potential. It is given as: U(r) = ∫ (G r
mM
)dr' = −G
2
+K .
mM
r
r'
Choosing the convention that the constant of integration K=0 assumes that the potential at infinity is defined to be 0.
The displacement of spring x is usually measured from the position of “neutral length ” or “relaxed length”. Often, it is most
convenient to identify this position as the origin of coordinate reference (x=0).
If the block is gently released from the stretched position (x = x ), energy conservation tells us that
f
kx = constant .
1 2 1 2 1 2
mv + kx =
2 2 2 f
If the block is released from the stretched position (x = x ), by the time the block reaches x=0 position, its speed will be
f
−
−
v(x = 0) = √
k
m
xf . The block will keep oscillating between x = −x and x .
f f
6.5.9 https://phys.libretexts.org/@go/page/14470
energy is: KE + PE = KE + PE . .
i i f f
If you know the potential energy for only some of the forces, then the conservation of energy law in its most general form must
be used: KE + PE + W + OE = KE + PE + OE
i i nc i f f f , where OE stands for all other energies.
Once you have solved a problem, always check the answer to see if it is reasonable.
Using the new energy conservation relationship
KEi + PEi + Wnc = KEf + PEf (6.5.13)
, we can apply the same problem-solving strategy as with the case of conservative forces.
The most important point is that the amount of nonconservative work equals the change in mechanical energy.
The work done by nonconservative (or dissipative) forces will irreversibly dissipated in the system.
Key Terms
potential: A curve describing the situation where the difference in the potential energies of an object in two different positions
depends only on those positions.
Coulomb force: the electrostatic force between two charges, as described by Coulomb’s law
potential: A curve describing the situation where the difference in the potential energies of an object in two different positions
depends only on those positions.
conservative force: A force with the property that the work done in moving a particle between two points is independent of the
path taken.
Hooke’s law: the principle that the stress applied to a solid is directly proportional to the strain produced. This law describes the
behavior of springs and solids stressed within their elastic limit.
conservation: A particular measurable property of an isolated physical system does not change as the system evolves.
isolated system: A system that does not interact with its surroundings, that is, its total energy and mass stay constant.
frictional force: Frictional force is the force resisting the relative motion of solid surfaces, fluid layers, and material elements
sliding against each other.
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
potential energy: The energy an object has because of its position (in a gravitational or electric field) or its condition (as a
stretched or compressed spring, as a chemical reactant, or by having rest mass)
conservative force: A force with the property that the work done in moving a particle between two points is independent of the
path taken.
dissipative force: A force resulting in dissipation, a process in which energy (internal, bulk flow kinetic, or system potential) is
transformed from some initial form to some irreversible final form.
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6.6: Power
What is Power?
In physics, power is the rate of doing work—the amount of energy consumed per unit time.
learning objectives
Relate power to the transfer, use, and transformation of different types of energy
In physics, power is the rate of doing work. It is the amount of energy consumed per unit time. The unit of power is the joule per
second (J/s), known as the watt (in honor of James Watt, the eighteenth-century developer of the steam engine). For example, the
rate at which a lightbulb transforms electrical energy into heat and light is measured in watts (W)—the more wattage, the more
power, or equivalently the more electrical energy is used per unit time.
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running because the work is done in a shorter amount of time. The output power of an electric motor is the product of the torque the
motor generates and the angular velocity of its output shaft. The power expended to move a vehicle is the product of the traction
force of the wheels and the velocity of the vehicle.
Examples of power are limited only by the imagination, because there are as many types as there are forms of work and energy.
Sunlight reaching Earth’s surface carries a maximum power of about 1.3 kilowatts per square meter (kW/m2). A tiny fraction of
this is retained by Earth over the long term. Our consumption rate of fossil fuels is far greater than the rate at which they are stored,
so it is inevitable that they will be depleted. Power implies that energy is transferred, perhaps changing form. It is never possible to
change one form completely into another without losing some of it as thermal energy. For example, a 60-W incandescent bulb
converts only 5 W of electrical power to light, with 55 W dissipating into thermal energy. Furthermore, the typical electric power
plant converts only 35 to 40 percent of its fuel into electricity. The remainder becomes a huge amount of thermal energy that must
be dispersed as heat transfer, as rapidly as it is created. A coal-fired power plant may produce 1,000 megawatts; 1 megawatt (MW)
is 106 W of electric power. But the power plant consumes chemical energy at a rate of about 2,500 MW, creating heat transfer to
the surroundings at a rate of 1,500 MW.
Coal-fired Power Plant: Tremendous amounts of electric power are generated by coal-fired power plants such as this one in
China, but an even larger amount of power goes into heat transfer to the surroundings. The large cooling towers here are needed to
transfer heat as rapidly as it is produced. The transfer of heat is not unique to coal plants but is an unavoidable consequence of
generating electric power from any fuel—nuclear, coal, oil, natural gas, or the like.
learning objectives
Identify what factors play a role in basal metabolic rate (BMR)
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Energy Conversion in Humans: Energy consumed by humans is converted to work, thermal energy, and stored fat. By far the
largest fraction goes to thermal energy, although the fraction varies depending on the type of physical activity.
Useful Work
Work done by a person is sometimes called useful work, which is work done on the outside world, such as lifting weights. Useful
work requires a force exerted through a distance on the outside world, and so it excludes internal work, such as that done by the
heart when pumping blood. Useful work does include that done in climbing stairs or accelerating to a full run, because these are
accomplished by exerting forces on the outside world. Forces exerted by the body are nonconservative, so that they can change the
mechanical energy (KE+PE) of the system worked upon, and this is often the goal.
For example, what is the power output for a 60.0-kg woman who runs up a 3.00 m high flight of stairs in 3.50 s, starting from rest
but having a final speed of 2.00 m/s?.
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Woman Running Up Stairs: When this woman runs upstairs starting from rest, she converts the chemical energy originally from
food into kinetic energy and gravitational potential energy. Her power output depends on how fast she does this.
Her power output depends on how fast she does this. The work going into mechanical energy is W= KE + PE. At the bottom of the
stairs, we take both KE and PEg as initially zero; thus,
1
2
W = KEf + PEg = mv + mgh (6.6.1)
f
2
where h is the vertical height of the stairs. Because all terms are given, we can calculate W and then divide it by time to get power.
Substituting the expression for W into the definition of power given in the previous equation, P = yields W
1 2
W mv + mgh
2 f
P = = (6.6.2)
t t
The woman does 1764 J of work to move up the stairs compared with only 120 J to increase her kinetic energy; thus, most of her
power output is required for climbing rather than accelerating.
Energy consumption is directly proportional to oxygen consumption because the digestive process is basically one of oxidizing
food. We can measure the energy people use during various activities by measuring their oxygen use. Approximately 20 kJ of
energy are produced for each liter of oxygen consumed, independent of the type of food.
Key Points
Power implies that energy is transferred, perhaps changing form.
Energy transfer can be used to do work, so power is also the rate at which this work is performed.
The unit of power is the joule per second (J/s), known as the watt.
The rate at which the body uses food energy to sustain life and to do different activities is called the metabolic rate, and the
corresponding rate when at rest is called the basal metabolic rate (BMR).
The energy included in the basal metabolic rate is divided among various systems in the body, with the largest fraction going to
the liver and spleen, and the brain coming next.
About 75% of food calories are used to sustain basic body functions included in the basal metabolic rate.
Work done by a person is sometimes called useful work, which is work done on the outside world, such as lifting weights.
The energy consumption of people during various activities can be determined by measuring their oxygen use, because the
digestive process is basically one of oxidizing food.
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Key Terms
power: A measure of the rate of doing work or transferring energy.
watt: In the International System of Units, the derived unit of power; the power of a system in which one joule of energy is
transferred per second.
basal metabolic rate: The amount of energy expended while at rest in a neutrally temperate environment, in the post-
absorptive state.
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6.7: CASE STUDY: World Energy Use
World Energy Use
The most prominent sources of energy used in the world are non-renewable (i.e., unsustainable).
learning objectives
Explain why renewable energy sources must be found and utilized
Energy Use
World energy consumption is the total amount of energy used by all humans on the planet (measured on a per-year basis). This
measurement is the sum of all energy sources (and purposes) in use. Who measures this? Several organizations publish this data,
including the International Energy Agency (IEA), the US Energy Information Administration (EIA), and the European
Environment Agency. This data is useful because evaluating this information to discover trends might yield energy issues not
currently being addressed, thereby encouraging the search for solutions. The IEA established a goal of limiting global warming to 2
degrees Celsius, but this goal is becoming more difficult to reach each year that the necessary action is not taken. In global energy
use, fossil fuels make up a substantial portion. In 2011 they received over $500 billion in subsidies—six times more than that
received by renewable energy sources.
Implementing new practices that will utilize different, renewable energy sources is important because having access to energy is
important—it maintains our quality of life. Fossil fuels, however, are not sustainable at the rate they are currently used. About 40%
of the world’s energy comes from oil, but oil prices are dependent on uncertain factors (such as availability, politics, and world
events). The United States alone uses 24% of the world’s oil per year, yet it makes up only 4.5% of the world’s population! In
2008, total worldwide energy consumption was 474 exajoules (474×1018 J=132,000 TWh)—equivalent to an average power usage
of 15 terawatts (1.504×1013 W). Potential renewable energy sources include: solar energy at 1600 EJ (444,000 TWh), wind power
at 600 EJ (167,000 TWh), geothermal energy at 500 EJ (139,000 TWh), biomass at 250 EJ (70,000 TWh), hydropower at 50 EJ
(14,000 TWh) and ocean energy at 1 EJ (280 TWh).
Types of Energy
shows a pie chart of world energy usage by category—both renewable and nonrenewable sources. Renewable energy comes from
sources with an unlimited supply. This includes energy from water, wind, the sun, and biomass. In the US, only 10% of energy
comes from renewable sources (mostly hydroelectric energy). Nonrenewable sources makes up 85% of worldwide energy usage—
from sources that eventually will be depleted, such as oil, natural gases and coal.
World Energy Use: This chart shows that the primary worldwide energy sources nonrenewable. If new practices are not put in
place now, this model will not be sustainable.
Energy Needs
In the last 50 years, the global energy demand has tripled due to the number of developing countries and innovations in technology.
It is projected to triple again over the next 30 years. In Europe, many in such developing areas recognize that the need for
renewable energy sources, as the present course of energy usage cannot be sustained indefinitely. While renewable energy
development makes up a only small percentage of the field, strides are being made in natural energy, particularly wind energy.
For example, by the year 2020 Germany plans to meet 10% of their total energy usage and 20% of its electricity usage with
renewable resources. While some countries are making improvements in this field, coal usage is still a huge problem. In China, two
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thirds of the energy used each year is from commercial coal energy. India imports 50% of its oil, and 70% of its electricity is
produced from coal, which is highly polluting.
Key Points
The energy consumption increases with the increasing number of developing areas. In order for this development to continue,
while maintaining quality of life, new and renewable energy sources must be found and utilized.
Renewable energy comes from sources that will never deplete, no matter how much is used. An example of this is wind energy,
which had been growing in popularity in countries like India and Germany.
Nonrenewable energy makes up 85% the energy used on earth—the most popular form of energy being oil.
Key Terms
fossil fuel: Any fuel derived from hydrocarbon deposits such as coal, petroleum, natural gas and, to some extent, peat; these
fuels are irreplaceable, and their burning generates the greenhouse gas carbon dioxide.
renewable energy: Energy that can be replenished at the same rate as it is used.
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6.8: Further Topics
learning objectives
Compare the different forms of energy interrelate to one another
6.8.1 https://phys.libretexts.org/@go/page/15489
Atomic bomb explosion: The mushroom cloud of the atomic bombing of Nagasaki, Japan
In each of the aforementioned forms, energy exists as either kinetic energy, potential energy, or a combination of both. It is
important to note that the above list is not necessarily complete as we may discover new forms of energy in the future such as “dark
energy. ” Also, each of the forms that energy can take on (as listed above) are not necessarily mutually exclusive. For example,
luminous energy is radiant energy.
6.8.2 https://phys.libretexts.org/@go/page/15489
Types of Energy: A brief overview of energy, kinetic energy, gravitational potential energy, and the work-energy theorem for
algebra-based physics students.
Energy Transformations
Energy transformation occurs when energy is changed from one form to another, and is a consequence of the first law of
thermodynamics.
learning objectives
Summarize the consequence of the first law of thermodynamics on the total energy of a system
Energy transformation occurs when energy is changed from one form to another. It is a consequence of the first law of
thermodynamics that the total energy of a given system can only be changed when energy is added or subtracted from the system.
Often it appears that energy has been lost from a system when it simply has been transformed. For example, an internal combustion
engine converts the potential chemical energy in gasoline and oxygen into heat energy. This heat energy is then converted to kinetic
energy, which is then used to propel the vehicle that is utilizing the engine. The technical term for a device that converts energy
from one form to another is a transducer.
Energy Transformation: These figures illustrate the concepts of energy loss and useful energy output.
When analyzing energy transformations, it is important to consider the efficiency of conversion. The efficiency of conversions
describes the ratio between the useful output and input of an energy conversion machine. Some energy transformations can occur
with an efficiency of essentially 100%. For example, imagine a pendulum in a vacuum. As illustrated in, when the pendulum’s
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mass reaches its maximum height, all if its energy exists in the form of potential energy. However, when the pendulum is at its
lowest point, all of its energy exists in the form of kinetic energy.
Pendulum: This animation shows the velocity and acceleration vectors for a pendulum. One may note that at the maximum height
of the pendulum’s mass, the velocity is zero. This corresponds to zero kinetic energy and thus all of the energy of the pendulum is
in the form of potential energy. When the pendulum’s mass is at its lowest point, all of its energy is in the form of kinetic energy
and we see its velocity vector has a maximum magnitude here.
Other energy transformations occur with a much lower efficiency of conversion. For example, the theoretical limit of the energy
efficiency of a wind turbine (converting the kinetic energy of the wind to mechanical energy) is 59%. The process of
photosynthesis is able to transform the light energy of the sun into chemical energy that can be used by a plant with an efficiency of
conversion of a mere 6%.
learning objectives
Derive the potential of a point charge
A potential energy curve plots the potential energy of an object as a function of that object’s position. For example, see. The system
under consideration is a closed system, so the total energy of the system remains constant. This means that the kinetic and potential
energy always have to sum to be the same value. We observe that the potential energy increases as the kinetic energy decreases and
vice versa. The utility of a potential energy curve is that we can quickly determine the potential energy of the object in question at a
given position.
Potential Energy Curve: This figure illustrates the potential energy of a particle as a function of position. The kinetic energy is
also shown and is abbreviated K.
Equipotential lines trace lines of equal potential energy. In, if you were to draw a straight horizontal line through the center, that
would be an equipotential line. In and, if you travel along an equipotential line, the electric potential will be constant.
6.8.4 https://phys.libretexts.org/@go/page/15489
Equipotential Lines for Two Equal and Opposite Point Charges: Electric field (blue) and equipotential lines (green) for two
equal and opposite charges
kQ
Let us examine the physical explanation for the equipotentials lines in. The equation for the potential of a point charge is V = , r
9 2 2
where V is the potential, k is a constant with a value of 8.99 · 10 N m /C , Q is the magnitude of the point charge, and ris the
distance from the charge. So, every point that is the same distance from the point charge will have the same electric potential
energy. Therefore, if we draw a circle around the point charge, every point on the circle will have the same potential energy.
Work (W) is a measure of the change in potential energy (ΔPE) : W = −\DeltaPE. Since the potential energy does not change
along an equipotential line, you do not need to do any work to move along one. However, you do need to do work to move from
one equipotential line to another. Recall that work is zero if force is perpendicular to motion; in the figures shown above, the forces
resulting from the electric field are in the same direction as the electric field itself. So we note that each of the equipotential lines
must be perpendicular to the electric field at every point.
Key Points
Thermal, chemical, electric, radiant, nuclear, magnetic, elastic, sound, mechanical, luminous, and mass are forms that energy
can exist in.
Energy exists as either kinetic energy, potential energy, or a combination of both.
We may discover new forms of energy (like “dark energy”) in the future.
The total energy of a given system can only be changed when energy is added or subtracted from the system.
Often it appears that energy has been lost from a system when it simply has been transformed.
The efficiency of conversions describes the ratio between the useful output of an energy conversion machine and the input.
A potential energy curve plots the potential energy of an object as a function of its position.
Equipotential lines trace lines of equal potential energy.
You do not need to do any work to move along an equipotential line.
Key Terms
fusion: A nuclear reaction in which nuclei combine to form more massive nuclei with the concomitant release of energy.
electromagnetic radiation: radiation (quantized as photons) consisting of oscillating electric and magnetic fields oriented
perpendicularly to each other, moving through space
fission: The process of splitting the nucleus of an atom into smaller particles; nuclear fission.
pendulum: A body suspended from a fixed support so that it swings freely back and forth under the influence of gravity; it is
commonly used to regulate various devices such as clocks.
first law of thermodynamics: A version of the law of conservation of energy, specialized for thermodynamical systems. It is
usually formulated by stating that the change in the internal energy of a closed system is equal to the amount of heat supplied to
the system, minus the amount of work done by the system on its surroundings.
potential energy: The energy an object has because of its position (in a gravitational or electric field) or its condition (as a
stretched or compressed spring, as a chemical reactant, or by having rest mass)
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
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CHAPTER OVERVIEW
This page titled 7: Linear Momentum and Collisions is shared under a not declared license and was authored, remixed, and/or curated by
Boundless.
1
7.1: Introduction
Linear Momentum
Linear momentum is the product of the mass and velocity of an object, it is conserved in elastic and inelastic collisions.
learning objectives
Calculate the momentum of two colliding objects
In classical mechanics, linear momentum, or simply momentum (SI unit kg m/s, or equivalently N s), is the product of the mass
and velocity of an object. Mathematically it is stated as:
p = mv (7.1.1)
(Note here that p and v are vectors. ) Like velocity, linear momentum is a vector quantity, possessing a direction as well as a
magnitude. Linear momentum is particularly important because it is a conserved quantity, meaning that in a closed system (without
any external forces ) its total linear momentum cannot change.
Because momentum has a direction, it can be used to predict the resulting direction of objects after they collide, as well as their
speeds. Momentum is conserved in both inelastic and elastic collisions. (Kinetic energy is not conserved in inelastic collisions but
is conserved in elastic collisions. ) It important to note that if the collision takes place on a surface with friction, or if there is air
resistance, we would need to account for the momentum of the bodies that would be transferred to the surface and/or air.
Let’s take a look at a simple, one-dimensional example: The momentum of a system of two particles is the sum of their momenta.
If two particles have masses m1 and m2, and velocities v1 and v2, the total momentum is:
p = p1 + p2 = m1 v1 + m2 v2 . (7.1.2)
Keep in mind that momentum and velocity are vectors. Therefore, in 1D, if two particles are moving in the same direction, v1 and
v2 have the same sign. If the particles are moving in opposite directions they will have opposite signs.
If two particles were moving on a plane we would choose our xy-plane to be on the plane of motion. We can then write the x and y
component of the total momentum as:
px = p1x + p2x = m1 v1x + m2 v2x (7.1.3)
If the 2D momentum vector is decomposed into two components, the equations for each component are reduced to its 1D
equivalents.
Momentum, like energy, is important because it is conserved. “Newton’s cradle” shown in is an example of conservation of
momentum. As we will discuss in the next concept (on Momentum, Force, and Newton’s Second Law ), in classical mechanics,
conservation of linear momentum is implied by Newton’s laws. Only a few physical quantities are conserved in nature. Studying
these quantities yields fundamental insight into how nature works.
Newton’s Cradle: Total momentum of the system (or Cradle) is conserved. (neglecting frictional loss in the system. )
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Momentum, Force, and Newton’s Second Law
In the most general form, Newton’s 2nd law can be written as F =
dp
dt
.
learning objectives
Relate Newton’s Second Law to momentum and force
In a closed system (one that does not exchange any matter with the outside and is not acted on by outside forces), the total
momentum is constant. This fact, known as the law of conservation of momentum, is implied by Newton’s laws of motion.
Suppose, for example, that two particles interact. Because of the third law, the forces between them are equal and opposite. If the
particles are numbered 1 and 2, the second law states that
dp1 dp2
=− (7.1.5)
dt dt
or
d
(p1 + p2 ) = 0 (7.1.6)
dt
Therefore, total momentum (p1+p2) is constant. If the velocities of the particles are u1 and u2 before the interaction, and afterwards
they are v1 and v2, then
m1 u1 + m2 u2 = m1 v1 + m2 v2 (7.1.7)
This law holds regardless of the nature of the interparticle (or internal) force, no matter how complicated the force is between
particles. Similarly, if there are several particles, the momentum exchanged between each pair of particles adds up to zero, so the
total change in momentum is zero.
where F net is the net external force, Δp is the change in momentum, and Δt is the change in time.
This statement of Newton’s second law of motion includes
familiarthe more
Fnet = ma\(asaspecialcase. W ecanderivethisf ormasf ollows. F irst, notethatthechangeinmomentum\(Δp is given
by Δp = Δ(mv) . If the mass of the system is constant, then Δ(mv) = mΔv . So for constant mass, Newton’s second law of
motion becomes
Δp mΔv
Fnet = = . (7.1.9)
Δt Δt
Because Δv
Δt
= a , we get the familiar equation F = ma when the mass of the system is constant. Newton’s second law of motion
net
stated in terms of momentum is more generally applicable because it can be applied to systems where the mass is changing, such as
rockets, as well as to systems of constant mass.
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Momentum in a Closed System: In a game of pool, the system of entire balls can be considered a closed system. Therefore, the
total momentum of the balls is conserved.
Impulse
Impulse, or change in momentum, equals the average net external force multiplied by the time this force acts.
learning objectives
Explain the relationship between change in momentum and the amount of time a force acts
Impulse
Forces produce either acceleration or deceleration on moving bodies, and the greater the force acting on an object, the greater its
change in velocity and, hence, the greater its change in momentum. However, changing momentum is also related to how long a
time the force acts. If a brief force is applied to a stalled automobile, a change in its momentum is produced. The same force
applied over an extended period of time produces a greater change in the automobile’s momentum. The quantity of impulse is force
× time interval, or in shorthand notation:
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Momentum & Impulse: A brief overview of momentum and impulse for high school physics students.
where F is the net force on the system, and Δt is the duration of the force.
From Newton’s 2nd law:
Δp
F = (Δp : change in momentum), (7.1.11)
Δt
change in momentum equals the average net external force multiplied by the time this force acts.
Δp = FΔt. (7.1.12)
looks like as a function of time for a ball bouncing off the floor. The area under the curve has units of momentum and is equal to
the impulse or change in momentum between times t and t . That area is equal to the area inside the rectangle bounded by
1 2
F , t , and t . Thus, the impulses and their effects are the same for both the actual and effective forces. Equivalently, we can
ef f 1 2
simple find the area under the curve F(t) between t and t to compute the impulse in mathematical form:
1 2
Force vs. Time: A graph of force versus time with time along the x-axis and force along the y-axis for an actual force and an
equivalent effective force. The areas under the two curves are equal.
t2
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Key Points
Like velocity, linear momentum is a vector quantity, possessing a direction as well as a magnitude.
Momentum, like energy, is important because it is a conserved quantity.
The momentum of a system of particles is the sum of their momenta. If two particles have masses m1 and m2, and velocities v1
and v2, the total momentum is p = p + p = m v + m v . .
1 2 1 1 2 2
In a closed system, without any external forces, the total momentum is constant.
The familiar equation F = ma is a special case of the more general form of the 2nd law when the mass of the system is
constant.
Momentum conservation holds (in the absence of external force) regardless of the nature of the interparticle (or internal) force,
no matter how complicated the force is between particles.
A small force applied for a long time can produce the same momentum change as a large force applied briefly, because it is the
product of the force and the time for which it is applied that is important.
A force produces an acceleration, and the greater the force acting on an object, the greater its change in velocity and, hence, the
greater its change in momentum. However, changing momentum is also related to how long a time the force acts.
In case of a time-varying force, impulse can be calculated by integrating the force over the time duration.
t2
Impulse = ∫ F(t)dt. .
t1
Key Terms
inelastic: (As referring to an inelastic collision, in contrast to an elastic collision. ) A collision in which kinetic energy is not
conserved.
elastic collision: An encounter between two bodies in which the total kinetic energy of the two bodies after the encounter is
equal to their total kinetic energy before the encounter. Elastic collisions occur only if there is no net conversion of kinetic
energy into other forms.
conservation: A particular measurable property of an isolated physical system does not change as the system evolves.
closed system: A physical system that doesn’t exchange any matter with its surroundings and isn’t subject to any force whose
source is external to the system.
momentum: (of a body in motion) the product of its mass and velocity.
impulse: The integral of force over time.
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7.2: Conservation of Momentum
Internal vs. External Forces
Net external forces (that are nonzero) change the total momentum of the system, while internal forces do not.
learning objectives
Contrast the effects of external and internal forces on linear momentum and collisions
entire system p should be constant in the absence of net external forces. Forces external to the system may change the total
tot
momentum when their sum is not 0, but internal forces, regardless of the nature of the forces, will not contribute to the change in
the total momentum. To analyze a mechanical system, it is important to recognize which forces are internal and which are external.
Once a mechanical system is clearly defined, it’s not hard to understand what part should be considered external.
Newton’s Cradle: Total momentum of the system (or Cradle) is conserved. (neglecting frictional loss in the system. )
External forces: forces caused by external agent outside of the system.
Internal forces: forces exchanged by the particles in the system.
To give you a better idea, let’s consider a simple example. We have two hockey pucks sliding across a frictionless surface, and we
neglect air resistance for simplicity. They collide with each other at t=0.
Let’s first list all the forces present in the system. There are mainly three kinds of forces: Gravity, normal force (between ice &
pucks), and frictional forces during the collision between the pucks
How should we define our system? In most cases, we would be interested in the motion of the pucks (and nothing else). Therefore,
our system consists of two pucks (and nothing else). All the rest of the universe becomes external. With this in mind, we can see
that gravity and normal forces are external, while the frictional forces between pucks are internal. Since all the external forces
cancel out with each other, there are no net external forces. (Gravity and normal force on each puck have the same magnitude, but
are in the opposite directions) Therefore, we conclude that the total momentum of the two pucks should be a conserved quantity.
In the previous example, it is worthwhile to note that we didn’t assume anything about the nature of the collision between the
two pucks. Without knowing anything about the internal forces (frictional forces during contact), we learned that the total
momentum of the system is a conserved quantity (p1 and p2 are momentum vectors of the pucks. ) In fact, this relation holds
true both in elastic or inelastic collisions. Whether the total kinetic energy of the pucks is conserved or not, total momentum is
conserved.
Also note that, in the previous example, if we include the rest of the Earth in our system, the gravity and normal forces
themselves become internal.
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Key Points
External forces are forces caused by external agent outside of the system.
Internal forces are forces exchanged by the objects in the system.
To determine what part should be considered external and internal, mechanical system should be clearly defined.
Key Terms
inelastic: (As referring to an inelastic collision, in contrast to an elastic collision. ) A collision in which kinetic energy is not
conserved.
elastic: referring to elastic collision, in contrast to inelastic collision. A collision in which kinetic energy is conserved
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7.3: Collisions
Conservation of Energy and Momentum
In an inelastic collision the total kinetic energy after the collision is not equal to the total kinetic energy before the collision.
learning objectives
Assess the conservation of total momentum in an inelastic collision
At this point we will expand our discussion of inelastic collisions in one dimension to inelastic collisions in multiple dimensions. It
is still true that the total kinetic energy after the collision is not equal to the total kinetic energy before the collision. While inelastic
collisions may not conserve total kinetic energy, they do conserve total momentum.
We will consider an example problem in which one mass (m ) slides over a frictionless surface into another initially stationary
1
mass (m ). Air resistance will be neglected. The following things are known:
2
m1 = 0.250kg, (7.3.1)
m2 = 0.400kg, (7.3.2)
v1 = 2.00m/s, (7.3.3)
v2 = 0m/s, (7.3.5)
where v is the initial velocity of the first mass, v ' is the final velocity of the first mass,
1 1 v2 is the initial velocity of the second
mass, and θ ' is the angle between the velocity vector of the first mass and the x-axis.
1
The object is to calculate the magnitude and direction of the velocity of the second mass. After this, we will calculate whether this
collision was inelastic or not.
Since there are no net forces at work (frictionless surface and negligible air resistance), there must be conservation of total
momentum for the two masses. Momentum is equal to the product of mass and velocity. The initially stationary mass contributes
no initial momentum. The components of velocities along the x-axis have the form v ⋅ cos θ , where θ is the angle between the
velocity vector of the mass of interest and the x-axis.
Expressing these things mathematically:
m1 v1 = m1 v1 ' ⋅ cos(θ1 ) + m2 v‘2 ⋅ cos(θ2 ). (Eq. 2) (7.3.7)
The components of velocities along the y-axis have the form v ⋅ sin θ , where θ is the angle between the velocity vector of the mass
of interest and the x-axis. By applying conservation of momentum in the y-direction we find:
0 = m1 v1 ' ⋅ sin(θ1 ) + m2 v‘2 ⋅ sin(θ2 ). (Eq. 3) (7.3.8)
Now let’ use Eq. 3 to solve for v' . Re-arranging Eq. 3, we find:
2
7.3.1 https://phys.libretexts.org/@go/page/14475
1 1
2 2
Initial Kinetic Energy = m1 ⋅ v + m2 ⋅ v = 0.5J. (7.3.11)
1 2
2 2
1 2
1 2
Final Kinetic Energy = m1 ⋅ v1 ' + m2 ⋅ v2 ' ≈ 0.43J. (7.3.12)
2 2
Since these values are not the same we know that it was an inelastic collision.
Collision Example: This illustrates the example problem in which one mass collides into another mass that is initially stationary.
Glancing Collisions
Glancing collision is a collision that takes place under a small angle, with the incident body being nearly parallel to the surface.
learning objectives
Identify necessary conditions for a “glancing collision”
A collision is short duration interaction between two bodies or more than two bodies simultaneously causing change in motion of
bodies involved due to internal forces acted between them during this. Collisions involve forces (there is a change in velocity ). The
magnitude of the velocity difference at impact is called the closing speed. All collisions conserve momentum. What distinguishes
different types of collisions is whether they also conserve kinetic energy. Line of impact – It is the line which is common normal
for surfaces are closest or in contact during impact. This is the line along which internal force of collision acts during impact and
Newton’s coefficient of restitution is defined only along this line.
When dealing with an incident body that is nearly parallel to a surface, it is sometimes more useful to refer to the angle between the
body and the surface, rather than that between the body and the surface normal (see ), in other words 90° minus the angle of
incidence. This small angle is called a glancing angle. Collision at glancing angle is called “glancing collision”.
7.3.2 https://phys.libretexts.org/@go/page/14475
Collision: Object is deflected after the collision withthe surface. The angles between the body and the surface normal areindicated
as α and β. The angles between the body and the surface are 90– α and 90– β.
Collisions can either be elastic, meaning they conserve both momentum and kinetic energy, or inelastic, meaning they conserve
momentum but not kinetic energy. An inelastic collision is sometimes also called a plastic collision.
A “perfectly-inelastic” collision (also called a “perfectly-plastic” collision) is a limiting case of inelastic collision in which the two
bodies stick together after impact.
The degree to which a collision is elastic or inelastic is quantified by the coefficient of restitution, a value that generally ranges
between zero and one. A perfectly elastic collision has a coefficient of restitution of one; a perfectly-inelastic collision has a
coefficient of restitution of zero.
learning objectives
Assess the relationship among the collision equations to derive elasticity
An elastic collision is a collision between two or more bodies in which the total kinetic energy of the bodies before the collision is
equal to the total kinetic energy of the bodies after the collision. An elastic collision will not occur if kinetic energy is converted
into other forms of energy. It important to understand how elastic collisions work, because atoms often undergo essentially elastic
collisions when they collide. On the other hand, molecules do not undergo elastic collisions when they collide. In this atom we will
review case of collision between two bodies.
The mathematics of an elastic collision is best demonstrated through an example. Consider a first particle with mass m and 1
velocity v and a second particle with mass m and velocity v . If these two particles collide, there must be conservation of
1i 2 2i
momentum before and after the collision. If we know that this is an elastic collision, there must be conservation of kinetic energy
by definition. Therefore, the velocities of particles 1 and 2 after the collision (v and v respectively) will be related to the initial
1f 2f
velocities by:
1
2
2
m1 ⋅ v
1i
+
1
2
2
m2 ⋅ v
2i
=
1
2
2
m1 ⋅ v
1f
+
1
2
2
m2 ⋅ v
2f
(due to conservation of kinetic energy)
and
m1 ⋅ v1i + m2 ⋅ v2i = m1 ⋅ v1f + m2 ⋅ v2f (due to conservation of momentum).
Since we have two equations, we are able to solve for any two unknown variables. In our case, we will solve for the final velocities
of the two particles.
By grouping like terms and canceling out the ½ terms, we can rewrite our conservation of kinetic energy equation as:
2 2 2 2
m1 ⋅ (v −v ) = m2 ⋅ (v −v ). (Eq.1) (7.3.13)
1i 1f 2f 2i
By grouping like terms from our conservation of momentum equation we can find:
If we then divide Eq. 1 by Eq. 2 and perform some cancelations we will find:
7.3.3 https://phys.libretexts.org/@go/page/14475
We can solve for v as:
1f
At this point we see that v is still an unknown variable. So we can fix this by plugging Eq. 4 into our initial conservation of
2f
momentum equation. Our conservation of momentum equation with Eq. 4 substituted in looks like:
m1 ⋅ v1i + m2 ⋅ v2i = m1 ⋅ (v2f + v2i − v1i ) + m2 ⋅ v2f . (Eq.5) (7.3.17)
At this point we have successfully solved for the final velocity of the second particle. We still need to solve for the velocity of the
first particle, so let us do that by plugging Eq. 6 into Eq. 4.
2 ⋅ m1 (m2 − m1 )
v1f = [ v1i + v2i ] + v2i − v1i . (Eq. 7) (7.3.19)
(m2 + m1 ) (m2 + m1 )
Elastic Collision of Two Unequal Masses: In this animation, two unequal masses collide and recoil.
learning objectives
Construct an equation for elastic collision
Overview
As stated previously, there is conservation of total kinetic energy before and after an elastic collision. If an elastic collision occurs
in two dimensions, the colliding masses can travel side to side after the collision (not just along the same line as in a one
dimensional collision). The general approach to solving a two dimensional elastic collision problem is to choose a coordinate
system in which the velocity components of the masses can be decomposed along perpendicular axes.
7.3.4 https://phys.libretexts.org/@go/page/14475
AP Physics - Collisions in Multiple Dimensions
Collisions in Multiple Dimensions: A brief introduction to problem solving of collisions in two dimensions using the law of
conservation of momentum.
Example 7.3.1:
In this example, we consider only point masses. These are structure-less particles that cannot spin or rotate. We will consider a
case in which no outside forces are acting on the system, meaning that momentum is conserved. We will consider a situation in
which one particle is initially at rest. This situation is illustrated in.
Illustration of Elastic Collision in Two Dimensions: In this illustration, we see the initial and final configurations of two
masses that undergo an elastic collision in two dimensions.
7.3.5 https://phys.libretexts.org/@go/page/14475
By defining the x-axis to be along the direction of the incoming particle, we save ourselves time in breaking that velocity
vector into its x- and y- components. Now let us consider conservation of momentum in the x direction:
p1x + p2x = p1x ' + p2x '(Eq. 1) (7.3.21)
In Eq. 1, the initial momentum of the incoming particle is represented by p , the initial momentum of the stationary particle is
1x
represented by p , the final momentum of the incoming particle is represented by p . and the final momentum of the initially
2x
′
1x
We can expand Eq. 1 by taking into account that momentum is equal to the product of mass and velocity. Also, we know that
p = 0 because the initial velocity of the stationary particle is 0.
2x
The components of velocities along the x-axis have the form v ⋅ cos θ , where θ is the angle between the velocity vector of the
particle of interest and the x-axis.
Therefore:
′
m1 v1 = m1 v1 ' ⋅ cos(θ1 ) + m2 v ⋅ cos(θ2 )(Eq. 2) (7.3.22)
2
The components of velocities along the y-axis have the form v ⋅ sin θ , where θ is the angle between the velocity vector of the
particle of interest (denoted in the following equations by subscript 1 or 2) and the x-axis. We can apply conservation of
momentum in the y-direction in a similar way to yield:
′ ′
0 = m1 v ⋅ sin(θ1 ) + m2 v ⋅ sin(θ2 )(Eq. 3) (7.3.23)
1 2
In finding Eq. 3, it was taken into consideration that the incoming particle had no component of velocity along the y-axis.
The general approach to finding the defining equations for an n-dimensional elastic collision problem is to apply conservation of
momentum in each of the n- dimensions. You can generate an additional equation by utilizing conservation of kinetic energy.
learning objectives
Distinguish examples of inelastic collision from elastic collisions
Overview
In an inelastic collision the total kinetic energy after the collision is not equal to the total kinetic energy before the collision. This is
in contrast to an elastic collision in which conservation of total kinetic energy applies. While inelastic collisions may not conserve
total kinetic energy, they do conserve total momentum.
Collisions
If two objects collide, there are many ways that kinetic energy can be transformed into other forms of energy. For example, in the
collision of macroscopic bodies, some kinetic energy is turned into vibrational energy of the constituent atoms. This causes a
heating effect and results in deformation of the bodies. Another example in which kinetic energy is transformed into another form
of energy is when the molecules of a gas or liquid collide. When this happens, kinetic energy is often exchanged between the
molecules’ translational motion and their internal degrees of freedom.
A perfectly inelastic collision happens when the maximum amount of kinetic energy in a system is lost. In such a collision, the
colliding particles stick together. The kinetic energy is used on the bonding energy of the two bodies.
7.3.6 https://phys.libretexts.org/@go/page/14475
Sliding Block Example
Let us consider an example of a two-body sliding block system. The first block slides into the second (initially stationary block). In
this perfectly inelastic collision, the first block bonds completely to the second block as shown. We assume that the surface over
which the blocks slide has no friction. We also assume that there is no air resistance. If the surface had friction or if there was air
resistance, one would have to account for the bodies’ momentum that would be transferred to the surface and/or air.
Inelastic Collision: In this animation, one mass collides into another initially stationary mass in a perfectly inelastic collision.
Writing about the equation for conservation of momentum, one finds:
ma ua + mb ub = (ma + mb )v (7.3.25)
where mais the mass of the incoming block, ua is the velocity of the incoming block, mbis the mass of the initially stationary block,
ubis the velocity of initially stationary block (0 m/s), and v is the final velocity the two body system. Solving for the final velocity,
ma ua + mb ub
v = . (7.3.26)
ma + mb
Taking into account that the blocks have the same mass and that the one of the blocks is initially stationary, the expression for the
final velocity of the system may be defined as:
ua
v = . (7.3.27)
2
learning objectives
Relate inelastic collision multiple dimension equations to the one dimension collisions you learned earlier
At this point we will expand our discussion of inelastic collisions in one dimension to inelastic collisions in multiple dimensions. It
is still true that the total kinetic energy after the collision is not equal to the total kinetic energy before the collision. While inelastic
collisions may not conserve total kinetic energy, they do conserve total momentum.
7.3.7 https://phys.libretexts.org/@go/page/14475
Example 7.3.2:
Examples of Collisions
We will consider an example problem, illustrated in, in which one mass (m1m1) slides over a frictionless surface into another
initially stationary mass (m2m2). Air resistance will be neglected. The following quantities are known:
Collision Example: This illustrates the example problem in which one mass collides into another mass that is initially
stationary.
m1 = 0.250kg, (7.3.28)
m2 = 0.400kg, (7.3.29)
v1 = 2.00m/s, (7.3.30)
v2 = 0m/s, (7.3.32)
where v1v1 is the initial velocity of the first mass, v′1v1′ is the final velocity of the first mass, v2v2is the initial velocity of the
second mass, and θ′1θ1′ is the angle between the velocity vector of the first mass and the x-axis.
The object is to calculate the magnitude and direction of the velocity of the second mass. After this, we will calculate whether
this collision was inelastic or not.
Since there are no net forces at work (frictionless surface and negligible air resistance), there must be conservation of total
momentum for the two masses. Momentum is equal to the product of mass and velocity. The (initially) stationary mass
7.3.8 https://phys.libretexts.org/@go/page/14475
contributes no initial momentum. The components of velocities along the x-axis have the form v⋅cosθv⋅cosθ, where θ is the
angle between the velocity vector of the mass of interest and the x-axis.
Expressing these things mathematically:
The components of velocities along the y-axis have the form v ⋅ sin θ , where θ is the angle between the velocity vector of the
mass of interest and the x-axis. By applying conservation of momentum in the y-direction we find:
Now let’ use Eq. 3 to solve for v' . Re-arranging Eq. 3, we find:
2
As these values are not the same, we know this was an inelastic collision.
Key Points
In an inelastic collision the total kinetic energy after the collision is not equal to the total kinetic energy before the collision.
If there are no net forces at work (collision takes place on a frictionless surface and there is negligible air resistance ), there
must be conservation of total momentum for the two masses.
The variable θ is the angle between the velocity vector of the mass of interest and the x-axis in traditional Cartesian coordinate
systems.
Collision is short duration interaction between two bodies or more than two bodies simultaneously causing change in motion of
bodies involved due to internal forces acted between them during this.
Collisions can either be elastic, meaning they conserve both momentum and kinetic energy, or inelastic, meaning they conserve
momentum but not kinetic energy.
When dealing with an incident body that is nearly parallel to a surface, it is sometimes more useful to refer to the angle between
the body and the surface, rather than that between the body and the surface normal.
An elastic collision will not occur if kinetic energy is converted into other forms of energy.
While molecules do not undergo elastic collisions, atoms often undergo elastic collisions when they collide.
If two particles are involved in an elastic collision, the velocity of the first particle after collision can be expressed as:
( m1 −m2 ) 2⋅m2
v1f = v1i + v2i .
( m2 +m1 ) ( m2 +m1 )
If two particles are involved in an elastic collision, the velocity of the second particle after collision can be expressed as:
2⋅m1 ( m2 −m1 )
v2f = v1i + v2i .
( m2 +m1 ) ( m2 +m1 )
If an elastic collision occurs in two dimensions, the colliding masses can travel side to side after the collision.
By defining the x- axis to be along the direction of the incoming particle, we can simplify the defining equations.
The general approach to finding the defining equations for an n-dimensional elastic collision problem is to apply conservation
of momentum in each of the n- dimensions. You can generate an additional equation by utilizing conservation of kinetic energy.
In an inelastic collision, the total kinetic energy after the collision is not equal to the total kinetic energy before the collision.
While inelastic collisions may not conserve total kinetic energy, they do conserve total momentum.
7.3.9 https://phys.libretexts.org/@go/page/14475
A perfectly inelastic collision happens when the maximum amount of kinetic energy in a system is lost.
In an inelastic collision the total kinetic energy after the collision is not equal to the total kinetic energy before the collision.
If there are no net forces at work (i.e., collision takes place on a frictionless surface and there is negligible air resistance ), there
must be conservation of total momentum for the two masses.
The variable θ is the angle between the velocity vector of the mass of interest and the x-axis in traditional Cartesian coordinate
systems.
Key Terms
kinetic energy: The energy possessed by an object because of its motion, equal to one half the mass of the body times the
square of its velocity.
momentum: (of a body in motion) the product of its mass and velocity.
force: A physical quantity that denotes ability to push, pull, twist or accelerate a body which is measured in a unit dimensioned
in mass × distance/time² (ML/T²): SI: newton (N); CGS: dyne (dyn)
elastic collision: An encounter between two bodies in which the total kinetic energy of the two bodies after the encounter is
equal to their total kinetic energy before the encounter. Elastic collisions occur only if there is no net conversion of kinetic
energy into other forms.
dimension: A measure of spatial extent in a particular direction, such as height, width or breadth, or depth.
degrees of freedom: A degree of freedom is an independent physical parameter, often called a dimension, in the formal
description of the state of a physical system. The set of all dimensions of a system is known as a phase space.
friction: A force that resists the relative motion or tendency to such motion of two bodies in contact.
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7.4: Rocket Propulsion
Rocket Propulsion, Changing Mass, and Momentum
In rocket propulsion, matter is forcefully ejected from a system, producing an equal and opposite reaction on what remains.
learning objectives
Identify physical principles of rocket propulsion
change in momentum of the system. ) So the center of mass of the system is in free fall but, by rapidly expelling mass, part of the
system can accelerate upward. It is a commonly held misconception that the rocket exhaust pushes on the ground. If we consider
thrust; that is, the force exerted on the rocket by the exhaust gases, then a rocket’s thrust is greater in outer space than in the
atmosphere or on the launch pad. In fact, gases are easier to expel into a vacuum.
Free-body diagram of rocket propulsion: (a) This rocket has a mass m and an upward velocity v. The net external force on the
system is −mg, if air resistance is neglected. (b) A time Δt later the system has two main parts, the ejected gas and the remainder of
the rocket. The reaction force on the rocket is what overcomes the gravitational force and accelerates it upward.
By calculating the change in momentum for the entire system over Δt, and equating this change to the impulse, the following
expression can be shown to be a good approximation for the acceleration of the rocket.
7.4.1 https://phys.libretexts.org/@go/page/14476
ve Δm
a = −g (7.4.1)
m Δt
where a is the acceleration of the rocket, ve is the escape velocity, m is the mass of the rocket, Δm is the mass of the ejected gas,
and Δt is the time in which the gas is ejected.
Factors of Acceleration
A rocket’s acceleration depends on three major factors, consistent with the equation for acceleration of a rocket. First, the greater
the exhaust velocity of the gases relative to the rocket, ve, the greater the acceleration is. The practical limit for ve is about
2.5×103m/s for conventional (non-nuclear) hot-gas propulsion systems. The second factor is the rate at which mass is ejected from
the rocket. This is the factor Δm
Δt
in the equation. The quantity ( Δm
Δt
)v , with units of newtons, is called “thrust.” The faster the
e
rocket burns its fuel, the greater its thrust, and the greater its acceleration. The third factor is the mass m of the rocket. The smaller
the mass is (all other factors being the same), the greater the acceleration. The rocket mass m decreases dramatically during flight
because most of the rocket is fuel to begin with, so that acceleration increases continuously, reaching a maximum just before the
fuel is exhausted.
To achieve the high speeds needed to hop continents, obtain orbit, or escape Earth’s gravity altogether, the mass of the rocket other
than fuel must be as small as possible. It can be shown that, in the absence of air resistance and neglecting gravity, the final velocity
of a one-stage rocket initially at rest is
m0
v = ve ln (7.4.2)
mr
where ln(m0/mr) is the natural logarithm of the ratio of the initial mass of the rocket (m0) to what is left (mr) after all of the fuel is
exhausted. (Note that v is actually the change in velocity, so the equation can be used for any segment of the flight. If we start from
rest, the change in velocity equals the final velocity. )
Key Points
The propulsion of all rockets is explained by the same physical principle: Newton’s third law of motion.
A rocket’s acceleration depends on three major factors: the exhaust velocity, the rate the exhaust is ejected, and the mass of the
rocket.
To achieve the high speeds needed to hop continents, obtain orbit, or escape Earth’s gravity altogether, the mass of the rocket
other than fuel must be as small as possible.
Key Terms
Newton’s third law of motion: states that all forces exist in pairs: if one object A exerts a force FA on a second object B, then
B simultaneously exerts a force FB on A, and the two forces are equal and opposite: FA = −FB .
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7.5: Center of Mass
learning objectives
Identify the center of mass for an object with continuous mass distribution
In the previous modules on “Center of Mass and Translational Motion,” we learned why the concept of center of mass (COM)
helps solving mechanics problems involving a rigid body. Here, we will study the rigorous definition of COM and how to
determine the location of it. The position of COM is mass weighted average of the positions of particles.
where rCOM and ri are vectors representing the position of COM and i-th particle respectively, and M and mi are the total mass
and mass of the i-th particle, respectively. This mean means that position of COM is mass weighted average of the positions of
particles.
where M is the total mass in the volume. If a continuous mass distribution has uniform density, which means ρ is constant, then the
center of mass is the same as the center of the volume.
Plumb Line Method for Center of Mass: Suspend the object from two locations and to drop plumb lines from the suspension points.
The intersection of the two lines is the center of mass.
In three dimensions: By supporting an object at three points and measuring the forces that resist the weight of the object, COM of
the three-dimensional coordinates of the center of mass can be determined.
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Motion of the Center of Mass
We can describe the translational motion of a rigid body as if it is a point particle with the total mass located at the COM—center of
mass.
learning objectives
Derive the center of mass for the translational motion of a rigid body
We can describe the translational motion of a rigid body as if it is a point particle with the total mass located at the center of mass
(COM). In this Atom. we will prove that the total mass (M) times the acceleration of the COM (aCOM), indeed, equals the sum of
external forces. That is,
You can see that the Newton’s 2nd law applies as if we are describing the motion of a point particle (with mass M) under the
influence of the external force.
Derivation
From the definition of the center of mass,
∑ mi ri
i
rCOM = , (7.5.4)
M
For example, when we confine our system to the Earth and the Moon, the gravitational force due to the Sun would be external,
while the gravitational force on the Earth due to the Moon (and vice versa) would be internal. Since the gravitational forces
between the Earth and the Moon are equal in magnitude and opposite in direction, they will cancel out each other in the sum (see ).
COM of the Earth and Moon: Earth and Moon orbiting a COM inside the Earth. The red cross represents the COM of the two-body
system. The COM will orbit around the Sun as if it is a point particle.
Corollary
When there is no external force, the COM momentum is conserved.
Proof: Since there is no external force,M ⋅ a COM =0 . Therefore,
M ⋅ vCOM = constant.
7.5.2 https://phys.libretexts.org/@go/page/14477
Proof
Since there is no external force,
M ⋅ aC OM = 0. (7.5.5)
Therefore,
M ⋅ vC OM = constant .
□
learning objectives
Estimate the COM of a given object
The center of mass (COM) is an important physical concept. It is the point on an object at which the weighted relative position of
the distributed mass sums to zero—the point about which objects rotate.
Human proportions have been important in art, measurement, and medicine (a well known drawing of the human body is seen in ).
Although the human body has complicated features, the location of the center of mass (COM) could be a good indicator of body
proportions. The center of mass of the human body depends on the gender and the position of the limbs. In a standing posture, it is
typically about 10 cm lower than the navel, near the top of the hip bones. In this Atom, we will learn how to measure the COM of a
human body.
Leonardo da Vinci’s “The Vitruvian Man”: Vitruvian Man: A drawing created by Leonardo da Vinci. The drawing is based on
the correlations of ideal human proportions with geometry described[4] by the ancient Roman architect Vitruvius in Book III of his
treatise De Architectura.
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Example 7.5.1:
First, let’s take two scales and a wooden beam (H meter long), long enough to contain the entire body of the subject. Put the
scales H meters apart, and place the beam across the scales, as illustrated in. Now, let the subject lie on the beam. Make sure
that his/her heels are aligned with one end of the beam. Measure the readings (F1, F2) on the scale.
The COM of a Human Body: This figure demonstrates measuring the COM of a human body.
The system (person+beam) has three external forces: gravity on the subject (FCM), and normal forces from the scales F1 and F2.
The equation of motion for force (F=ma) will give us:
F1 + F2 = Mg, (7.5.6)
where M is mass of the subject. (We assume that the wooden beam has no mass. ) This equation doesn’t provide all the
information to locate the COM. However, the equation of motion for torque (τ = Iα) helps.
Since the net torque of the system is zero (hence no rotational acceleration),
hF2 − (H − h)F1 = 0.(h : COM height) (7.5.7)
The COM is chosen as the origin for the torque. Therefore, gravity contributes nothing as a torque. Solving for h and using the
equation of motion for force, we get
HF1
h= . (7.5.8)
Mg
learning objectives
Support the presence of COM in three dimensional bodies in motion
7.5.4 https://phys.libretexts.org/@go/page/14477
Forces on the COM: Left: The force appears to operate on the COM is “mgsinθ. Right: The force appears to operate on the COM
is “mg”.
Describing motions of parts or particles that have different motions would be quite complicated to do in an integrated manner.
However, such three dimensional bodies in motion have one surprising, simplifying characteristic—a geometric point that behaves
like a particle. This point is known as center of mass, abbreviated COM (the mathematical definition of COM will be introduced in
the next Atom on “Locating the Center of Mass”). It has the following two characterizing aspects:
The center of mass appears to carry the whole mass of the body.
At the center of mass, all external forces appear to apply.
Significantly, the center of a ball (the COM of a rolling ball) follows a straight linear path; whereas the COM of a stick follows a
parabolic path (as shown in the figure above). Secondly, the forces appear to operate on the COMs in two cases (“mgsinθ and
“mg”) as if they were indeed particle-like objects. This concept of COM, therefore, eliminate the complexities otherwise present in
attempting to describe motions of rigid bodies.
Key Points
The center of mass (COM) is a statement of spatial arrangement of mass (i.e. distribution of mass within the system).
The experimental determination of the center of mass of a body uses gravity forces on the body and relies on the fact that in the
parallel gravity field near the surface of the earth the center of mass is the same as the center of gravity.
For a 2D object, an experimental method for locating the center of mass is to suspend the object from two locations and to drop
plumb lines from the suspension points. The intersection of the two lines is the center of mass.
The total mass times the acceleration of the center of mass equals the sum of external forces.
For the translational motion of a rigid body with mass M, Newton’s 2nd law applies as if we are describing the motion of a
point particle (with mass M) under the influence of the external force.
When there is no external force, the center of mass momentum is conserved.
Although a human body has complicated features, the location of the center of mass (COM) could be a good indicator of the
body proportions.
We can measure the location of COM with two scales and a wooden beam. The linear and rotational equations of motion gives
us the location.
The center of mass of the human body depends on the gender and the position of the limbs. In a standing posture, it is typically
about 10 cm lower than the navel, near the top of the hip bones.
In a motion of a rigid body, different parts of the body have different motions.This means that these bodies may not behave like
a point particle.
There is a characteristic geometric point of the three dimensional body in motion. This point behaves as a particle, and is known
as center of mass, abbreviated COM. COM appears to carry the whole mass of the body. All external forces appear to apply at
COM.
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To describe the motion of a rigid body (with possibly a complicated geometry), we separate the translational part of the motion
from the rotational part.
Key Terms
plumb line: A cord with a weight attached, used to produce a vertical line.
rigid body: An idealized solid whose size and shape are fixed and remain unaltered when forces are applied; used in Newtonian
mechanics to model real objects.
center of mass: The center of mass (COM) is the unique point at the center of a distribution of mass in space that has the
property that the weighted position vectors relative to this point sum to zero.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
point particle: An idealization of particles heavily used in physics. Its defining feature is that it lacks spatial extension,
meaning that geometrically the particle is equivalent to a point.
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CHAPTER OVERVIEW
This page titled 8: Static Equilibrium, Elasticity, and Torque is shared under a not declared license and was authored, remixed, and/or curated by
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1
8.1: Introduction
learning objectives
Describe the effect of the torque on an object
Torque about a point is a concept that denotes the tendency of force to turn or rotate an object in motion. This tendency is measured
in general about a point, and is termed as moment of force. The torque in angular motion corresponds to force in translation. It is the
“cause” whose effect is either angular acceleration or angular deceleration of a particle in general motion. Quantitatively, it is
defined as a vector given by:
AP Physics 1 - Torque
Torque: A brief introduction to torque for students studying rotational motion in algebra-based physics courses such as AP Physics
1 and Honors Physics.
mathrmT = r × F (8.1.1)
Rotation is a special case of angular motion. In the case of rotation, torque is defined with respect to an axis such that vector “r” is
constrained as perpendicular to the axis of rotation. In other words, the plane of motion is perpendicular to the axis of rotation.
Clearly, the torque in rotation corresponds to force in translation.
Torque is the cross product of force cross length of the moment arm; it is involved whenever there is a rotating object. Torque can
also be expressed in terms of the angular acceleration of the object.
The determination of torque’s direction is relatively easier than that of angular velocity. The reason for this is simple: the torque
itself is equal to vector product of two vectors, unlike angular velocity which is one of the two operands of the vector product.
8.1.1 https://phys.libretexts.org/@go/page/14480
Clearly, if we know the directions of two operands here, the direction of torque can easily be interpreted.
Since torque depends on both the force and the distance from the axis of rotation, the SI units of torque are newton-meters.
Key Points
Torque is found by multiplying the applied force by the distance to the axis of rotation, called the moment arm.
Torque is to rotation as force is to motion.
The unit of torque is the newton-meter.
Key Terms
vector: A directed quantity, one with both magnitude and direction; the between two points.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
angular motion: The motion of a body about a fixed point or fixed axis (as of a planet or pendulum). It is equal to the angle
passed over at the point or axis by a line drawn to the body.
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8.2: Conditions for Equilibrium
learning objectives
Identify the first condition of equilibrium
In order to achieve this conditon, the forces acting along each axis of motion must sum to zero. For example, the net external forces
along the typical x– and y-axes are zero. This is written as
net Fx = 0 and net F y =0 .
The condition F = 0 must be true for both static equilibrium, where the object’s velocity is zero, and dynamic equilibrium,
net
8.2.1 https://phys.libretexts.org/@go/page/14481
A Car in Dynamic Equilibrium: This car is in dynamic equilibrium because it is moving at constant velocity. The forces in all
directions are balanced.
Second Condition
The second condition of static equilibrium says that the net torque acting on the object must be zero.
learning objectives
Identify the second condition of static equilibrium
A child’s seesaw, shown in, is an example of static equilibrium. An object in static equilibrium is one that has no acceleration in
any direction. While there might be motion, such motion is constant.
Two children on a seesaw: The system is in static equilibrium, showing no acceleration in any direction.
If a given object is in static equilibrium, both the net force and the net torque on the object must be zero. Let’s break this down:
Two-Component Forces
In equilibrium, the net force and torque in any particular direction equal zero.
8.2.2 https://phys.libretexts.org/@go/page/14481
learning objectives
Calculate the net force and the net torque for an object in equilibrium
An object with constant velocity has zero acceleration. A motionless object still has constant (zero) velocity, so motionless objects
also have zero acceleration. Newton’s second law states that:
∑ F = ma (8.2.1)
so objects with constant velocity also have zero net external force. This means that all the forces acting on the object are balanced
— that is to say, they are in equilibrium.
This rule also applies to motion in a specific direction. Consider an object moving along the x-axis. If no net force is applied to the
object along the x-axis, it will continue to move along the x-axis at a constant velocity, with no acceleration.
Car Moving at Constant Velocity: A moving car for which the net x and y force components are zero
We can easily extend this rule to the y-axis. In any system, unless the applied forces cancel each other out (i.e., the resultant force is
zero), there will be acceleration in the direction of the resultant force. In static systems, in which motion does not occur, the sum of
the forces in all directions always equals zero. This concept can be represented mathematically with the following equations:
∑ Fx = max = 0 (8.2.2)
∑ Fy = may = 0 (8.2.3)
This rule also applies to rotational motion. If the resultant moment about a particular axis is zero, the object will have no rotational
acceleration about the axis. If the object is not spinning, it will not start to spin. If the object is spinning, it will continue to spin at
the same constant angular velocity. Again, we can extend this to moments about the y-axis as well. We can represent this rule
mathematically with the following equations:
∑ τx = Iαx = 0 (8.2.4)
∑ τy = Iαy = 0 (8.2.5)
Key Points
There are two conditions that must be met for an object to be in equilibrium.
The first condition is that the net force on the object must be zero for the object to be in equilibrium.
If net force is zero, then net force along any direction is zero.
The second condition necessary to achieve equilibrium involves avoiding accelerated rotation.
A rotating body or system can be in equilibrium if its rate of rotation is constant and remains unchanged by the forces acting on
it.
The magnitude of torque about a axis of rotation is defined to be τ = rF sin θ .
In equilibrium, the net force in all directions is zero.
If the net moment of inertia about an axis is zero, the object will have no rotational acceleration about the axis.
In each direction, the net force takes the form: ∑ F = ma = 0 and the net torque take the form: ∑ τ = Iα = 0 where the sum
represents the vector sum of all forces and torques acting.
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Key Terms
force: A physical quantity that denotes ability to push, pull, twist or accelerate a body which is measured in a unit dimensioned
in mass × distance/time² (ML/T²): SI: newton (N); CGS: dyne (dyn)
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
translation: Motion of a body on a linear path, without deformation or rotation, i.e. such that every part of the body moves at
the same speed and in the same direction; also (in physics), the linear motion of a body considered independently of its rotation.
equilibrium: The state of a body at rest or in uniform motion, the resultant of all forces on which is zero.
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8.3: Stability
learning objectives
Explain the relationship between how center of mass is defined and static equilibrium
For an object to be in static equilibrium, we expect it to stay in the same state indefinitely. If it starts accelerating away from its
current position, it would hardly be in equilibrium. To quantify equilibrium for a single object, there are two conditions:
1. The net external force on the object is zero: ∑ F = F = 0
i i net
Those two conditions hold regardless of whether the object we are talking about is a single point particle, a rigid body, or a
collection of discrete particles. Being in equilibrium means that we expect no changes to the linear momentum or the angular
momentum. Note that this does not mean that the system is not moving or rotating; instead it simply means that its movement will
not change as time goes on.
In a special case when the external forces are governed by some potential (e.g. gravitational potential) we can gain insight into the
dU(x)
nature of the equilibrium. From the definition of a potential we know that F = − | . When the first derivative is zero, we
ext
dx x0
can take the second derivative to find whether the equilibrium is stable or unstable. Explicitly, if the potential is concave-up at x0,
2
d U(x)
dx
2
|
x0
, then the system is stable; conversely, if the potential is concave-down, then the equilibrium is unstable. If the second
>0
derivative is zero or does not exist, then the equilibrium is neutral—neither stable nor unstable.
Mathematically, we can view this as a Taylor series expansion on the force slightly away from equilibrium,
2
dF(x) dU(x) d U(x)
F(x0 + δx) = F(x0 ) +
dx
∣x0 δxF(x0 + δx) = −
dx
∣x0 +(−
dx
2
∣x0 )δx ,
and when it is initially at equilibrium,
2
dU(x) d U(x)
F(x0 ) = 0F(x0 ) = 0F(x0 + δx) = −
dx
∣x
0
+(−
dx
2
∣x )δxU(y) = mgy
0
.
If the ball is at the top of the hill (where the potential is concave-down) it is possible for it to be perfectly balanced, and therefore at
equilibrium. But if it gets pushed just slightly to the side, then it will roll down the hill with increasing speed, and the equilibrium is
unstable.
Unstable Equilibrium: A ball on top of a hill can initially be balanced, but if it moves slightly left or right, it gets pushed further
and further away from the initial equilibrium position. This is an example of unstable equilibrium.
Our notion of “balance” comes directly from the formulation of equilibrium. For something to be “balanced” means that the net
external forces are zero. For example, a coin could balance standing up on a table. Initially the coin will feel no net external force
or torque; it is in equilibrium. But if pushed slightly to the side, it will become “off-balance,” experiencing both a force and a
torque causing it to fall to the table. It might have been initially “balanced” and at equilibrium, but it was an unstable equilibrium,
prone to being disturbed. But why all this talk of external forces, with no mention of internal forces? The reason is that all the
internal forces must sum to zero. This follows directly from Newton’s Third Law, F = −F . Every time we consider a force
12 21
from particle 1 on particle 2 inside of a system, we know that it will later be cancelled out by the corresponding force from particle
2 on particle 1. We could include those forces in the sum, but it is unnecessary and internal forces are often more complicated than
internal forces.
This differentiation between internal and external forces is a powerful one. It also implies that you can trace the motion of the
system as a whole (ignoring motion inside the system) through the net external force acting on a center of mass. A center of mass
acts as if it has the entire mass of the system, located at one point, and only feels external forces. Its position is defined as the
8.3.1 https://phys.libretexts.org/@go/page/14482
R=∑i mi ri
weighted average of all the particles in the system: ∑ mi
or if we have a continuous density of mass, , then we can
ρ(r)
i
∫ Vρ(r)rdV
integrate: R = ∫ Vρ(r)dV
. The power of the center of mass is that it hides all the details of what is happening internally. We do not
always want to lose the information of what is happening internally, but it is a useful tool to remember, when dealing with a
number of complicated interactions.
Key Points
Equilibrium is defined by no net forces or torques.
Stability of an equilibrium can be determined by the second derivative of the potential.
Defining a center of mass allows a simple way to study the behavior of a system or object as a whole.
Stable equilibrium requires a restoring force. This restoring force can be derived by a Taylor expansion of the force, F(x).
Key Terms
stable equilibrium: The response [of a system in static equilibrium] to a small perturbation is forces that tend to restore the
equilibrium.
center of mass: The center of mass (COM) is the unique point at the center of a distribution of mass in space that has the
property that the weighted position vectors relative to this point sum to zero.
static equilibrium: the physical state in which all components of a system are at rest and the net force is equal to zero
throughout the system
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8.4: Solving Statics Problems
learning objectives
Formulate and apply six steps to solve static problems
Statics is the study of forces in equilibrium. Recall that Newton’s second law states:
∑ F = ma (8.4.1)
Therefore, for all objects moving at constant velocity (including a velocity of 0 — stationary objects), the net external force is zero.
There are forces acting, but they are balanced — that is to say, they are “in equilibrium. ”
When solving equilibrium problems, it might help to use the following steps:
1. First, ensure that the problem you’re solving is in fact a static problem—i.e., that no acceleration (including angular
acceleration) is involved Remember:∑ F = ma = 0 for these situations. If rotational motion is involved, the condition
∑ τ = Iα = 0 must also be satisfied, where is torque, is the moment of inertia, and is the angular acceleration.
2. Choose a pivot point. Often this is obvious because the problem involves a hinge or a fixed point. If the choice is not obvious,
pick the pivot point as the location at which you have the most unknowns. This simplifies things because forces at the pivot
point create no torque because of the cross product:τ = rF
3. Write an equation for the sum of torques, and then write equations for the sums of forces in the x and y directions. Set these
sums equal to 0. Be careful with your signs.
4. Solve for your unknowns.
5. Insert numbers to find the final answer.
6. Check if the solution is reasonable by examining the magnitude, direction, and units of the answer. The importance of this last
step cannot be overstated, although in unfamiliar applications, it can be more difficult to judge reasonableness. However, these
judgments become progressively easier with experience.
Key Points
First, ensure that the problem you’re solving is in fact a static problem—i.e., that no acceleration (including angular
acceleration ) is involved.
Choose a pivot point — use the location at which you have the most unknowns.
Write equations for the sums of torques and forces in the x and y directions.
Solve the equations for your unknowns algebraically, and insert numbers to find final answers.
Key Terms
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
moment of inertia: A measure of a body’s resistance to a change in its angular rotation velocity
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8.5: Applications of Statics
learning objectives
Develop an understanding of how a machine applies force to work against a load force
Simple Machines
A simple machine is a device that changes the direction or magnitude of a force. They can be described as the simplest mechanisms
that use mechanical advantage (or leverage) to multiply force. Usually, the term “simple machine” is referring to one of the six
classical simple machines, defined by Renaissance scientists.
Simple Machines, Pulleys; Forces, Distances and MA: Describes the following terms as they relate to simple machine; input
force, output force, input distance, output distance, mechanical advantage.
Simple machines are devices used to multiply or augment a force that we apply—often at the expense of a distance through which
we apply the force. Some common examples include:
Lever
8.5.1 https://phys.libretexts.org/@go/page/14484
Wheel and Axle
Pulley
Inclined Plane
Wedge
Screw
When a device with a specific movement, called a mechanism, is joined with others to form a machine, these machines can be
broken down into elementary movements. For example, a bicycle is a mechanism made up of wheels, levers, and pulleys.
Simple Machines: Table of simple mechanisms, from Chambers’ Cyclopedia, 1728. [1] Simple machines provide a “vocabulary”
for understanding more complex machines.
Mechanics
A simple machine has an applied force that works against a load force. If there are no frictional losses, the work done on the load is
equal to the work done by the applied force. This allows an increase in the output force at the cost of a proportional decrease in
distance moved by the load. The ratio of the output force to the input force is the mechanical advantage of the machine. If the
machine does not absorb energy, its mechanical advantage can be calculated from the machine’s geometry. For instance, the
mechanical advantage of a lever is equal to the ratio of its lever arms.
Simple machines which do not experience frictional losses are called ideal machines. For these ideal machines, the power in (rate
of energy input) in equal to the power out (rate of energy output): Pin=PoutPin=Pout.
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Lever: The amount of force produced by a machine can not be greater than the amount of force put into it.
Further Examples
Wheelbarrows and shovels are also examples of simple machines (these utilize levers). They use only three forces: the input force,
output force, and force on the pivot. In the case of wheelbarrows, the output force is between the pivot (wheel’s axle) and the input
force. In the shovel, the input force is between the pivot and the load.
Examples of Simple Machines: Both of these machines use the concept of levers.
8.5.3 https://phys.libretexts.org/@go/page/14484
Arches and Domes
Arches and domes are structures that exhibit structural strength and can span large areas with no intermediate supports.
learning objectives
Explain how an arch exhibits structural strength and how a dome can span a large area without intermediate supports
Arches and domes are structures that exhibit structural strength and can span large areas with no intermediate supports. In this
atom, we will discuss the history and physics behind arches and domes.
Arches
An arch is a structure that spans a space, and supports structure and weight above it. Arches have been being built from as long ago
as the second millennium, but were not used for a variety of structures until the Romans took advantage of their capabilities.
Arches are a pure compression form. They span large areas by resolving forces into compressive stresses and eliminating tensile
stresses (referred to as arch action). As the forces in an arch are carried toward the ground, the arch will push outward at the base
(called thrust ). As the height of the arch decreases, the outward thrust increases. To prevent the arch from collapsing, the thrust
needs to be restrained, either with internal ties or external bracing. This external bracing is often called an abutment, as shown in.
Arches: A masonry arch1. Keystone 2. Voussoir 3. Extrados 4. Impost 5. Intrados 6. Rise 7. Clear span 8. Abutment
The most common true arch configurations are the fixed arch, the two-hinged arch and the three-hinged arch. The fixed arch is
most often used in reinforced concrete bridge and tunnel construction, where the spans are short. Because it is subject to additional
internal stress caused by thermal expansion and contraction, this type of arch is considered to be statically indeterminate. The two-
hinged arch is most often used to bridge long spans. This type of arch has pinned connections at the base. Unlike the fixed arch, the
pinned base is able to rotate, allowing the structure to move freely and compensate for the thermal expansion and contraction
caused by changes in outdoor temperature. Because the structure is pinned between the two base connections, which can result in
additional stresses, the two-hinged arch is also statically indeterminate, although not to the degree of the fixed arch.
Domes
A dome is an element of architecture that resembles the hollow upper half of a sphere. Dome structures made of various materials
(from mud to stone, wood, brick, concrete, metal, glass and plastic) and have a long architectural lineage extending into prehistory.
A dome is basically an arch that has been rotated around its central vertical axis. Domes have the same properties and capabilities
of arches, they can span large areas without intermediate supports and have a great deal of structural strength. When the base of a
dome is not the same shape as its supporting walls, for example when a circular dome is on a square structure, techniques are
employed to transition between the two. Pendentives are triangular sections of a sphere used to transition from the flat surfaces of
supporting walls to the round base of a dome.
Domes can be divided into two kinds, simple and compound. Simple domes use pendentives that are part of the same sphere as the
dome itself. Compound domes are part of the structure of a large sphere below that of the dome itself, forming a circular base, as
shown in.
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Compound Dome: A compound dome (red) with pendentives (yellow) from a sphere of greater radius than the dome.
learning objectives
Explain the forces exerted by muscles
8.5.5 https://phys.libretexts.org/@go/page/14484
The Forearm of a Person Holding a Book: (a. ) The biceps exert a force FB to support the weight of the forearm and the book.
The triceps are assumed to be relaxed. (b.) An approximately equivalent mechanical system with the pivot at the elbow joint
Very large forces are also created in the joints. Because muscles can contract but not expand beyond their resting length, joints and
muscles often exert forces that act in opposite directions, and thus subtract. Forces in muscles and joints are largest when their load
is far from the joint. For example, in racquet sports like tennis, the constant extension of the arm during game play creates large
forces. The mass times the lever arm of a tennis racquet is an important factor, and many players use the heaviest racquet they can
handle. It is no wonder that joint deterioration and damage to the tendons in the elbow, such as ‘tennis elbow,’ can result from
repetitive motion, undue torques, and possible poor racquet selection in such sports.
Various tried techniques for holding and using a racquet, bat, or stick can not only increase sporting prowess but can minimize
fatigue and long-term damage to the body. Training coaches and physical therapists use the knowledge of the relationships between
forces and torques in the treatment of muscles and joints. In physical therapy, an exercise routine can apply a particular force and
torque, which can revive muscles and joints in time. Some exercises should be performed under water, thus requiring the exertion
of more force and further strengthening muscles.
Key Points
The six classifications of simple machines were established by renaissance scientists; they are as follows: lever, wheel and axle,
pulley, inclined plane, wedge and screw.
Simple machines can be joined with other devices to create a more complicated machine. These building blocks are used to
explain how machines work.
The force output by a simple machine can exceed the force that was put into the machine.
Arches span large areas by resolving forces into compressive stresses and eliminating tensile stresses.
The most common true arch configurations are the fixed arch, the two-hinged arch, and the three-hinged arch.
A dome is basically an arch that has been rotated around its central vertical axis.
Domes are basically arches that have been rotated on their vertical axis, and have the same capabilities and properties of arches.
Domes can be divided into two kinds, simple and compound.
It is helpful to view muscles as a simple machines and draw them as free body diagrams.
In muscles, the input force is often much greater than the output force.
Very large forces are also created in the joints. Because muscles can contract but not expand beyond their resting length, joints
and muscles often exert forces that act in opposite directions, and thus subtract.
Key Terms
machine: A mechanical or electrical device that performs or assists in the performance of human tasks, whether physical or
computational, laborious or for entertainment.
leverage: A force amplified by means of a lever rotating around a pivot.
mechanical advantage: In a simple machine, the ratio of the output force to the input force.
compressive stress: Stress on materials that leads to a smaller volume.
tensile stress: Stress state leading to expansion; that is, the length of a material tends to increase in the tensile direction while
the volume remains constant.
pendentive: The concave triangular sections of vaulting that provide the transition between a dome and the square base on
which it is set and transfer the weight of the dome.
muscle: A contractile form of tissue which animals use to effect movement.
joint: Any part of the body where two bones join, in most cases allowing that part of the body to be bent or straightened.
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8.6: Elasticity, Stress, Strain, and Fracture
learning objectives
Identify properties of elastic objects
We now move from consideration of forces that affect the motion of an object (such as friction and drag) to those that affect an
object’s shape. If a bulldozer pushes a car into a wall, the car will not move once it hits the wall, but it will noticeably change
shape. A change in shape due to the application of a force is a deformation. Even very small forces are known to cause some
deformation. For small deformations, two important characteristics are observed. First, the object returns to its original shape when
the force is removed—that is, the deformation is elastic for small deformations. Second, the size of the deformation is proportional
to the force—that is, for small deformations, Hooke’s law is obeyed. In equation form, Hooke’s law is given by F = kΔL , where
ΔL is the change in length.
Elasticity is a measure of how difficult it is to stretch an object. In other words it is a measure of how small kk is. Very elastic
materials like rubber have small kk and thus will stretch a lot with only a small force.
Stress is a measure of the force put on the object over the area.
Strain is the change in length divided by the original length of the object.
Experiments have shown that the change in length (ΔL) depends on only a few variables. As already noted, ΔL is proportional to
the force F and depends on the substance from which the object is made. Additionally, the change in length is proportional to the
original length L0 and inversely proportional to the cross-sectional area of the wire or rod. For example, a long guitar string will
stretch more than a short one, and a thick string will stretch less than a thin one.
Tension/Compression: Tension: The rod is stretched a length ΔL when a force is applied parallel to its length. (b) Compression:
The same rod is compressed by forces with the same magnitude in the opposite direction. For very small deformations and uniform
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materials, ΔL is approximately the same for the same magnitude of tension or compression. For larger deformations, the cross-
sectional area changes as the rod is compressed or stretched.
Fracture
Fracture is caused by a strain placed on an object such that it deforms beyond its elastic limit and breaks.
learning objectives
Relate fracture with the elastic limit of a material
Materials cannot stretch forever. When a strain is applied to a material it deforms elastically proportional to the force applied.
However, after it has deformed a certain amount, the object can no longer take the strain and will break or fracture. The zone in
which it bends under strain is called the elastic region. In that region the object will bend and then return to its original shape when
the force is abated. Past that point, if more strain is added, the object may permanently deform and eventually fracture.
Fracture strength, also known as breaking strength, is the stress at which a specimen fails via fracture. This is usually determined
for a given specimen by a tensile test, which charts the stress-strain curve. The final recorded point is the fracture strength.
Fracture: This is a graph of deformation ΔL versus applied force F. The straight segment is the linear region where Hooke’s law is
obeyed. The slope of the straight region is 1k. For larger forces, the graph is curved but the deformation is still elastic—L will
return to zero if the force is removed. Still greater forces permanently deform the object until it finally fractures. The shape of the
curve near fracture depends on several factors, including how the force F is applied. Note that in this graph the slope increases just
before fracture, indicating that a small increase in F is producing a large increase in L near the fracture.
Bones, on the whole, do not fracture due to tension or compression. Rather they generally fracture due to sideways impact or
bending, resulting in the bone shearing or snapping. The behavior of bones under tension and compression is important because it
determines the load the bones can carry. Bones are classified as weight-bearing structures such as columns in buildings and trees.
Weight-bearing structures have special features; columns in building have steel-reinforcing rods while trees and bones are fibrous.
The bones in different parts of the body serve different structural functions and are prone to different stresses. Thus, the bone in the
top of the femur is arranged in thin sheets separated by marrow while, in other places, the bones can be cylindrical and filled with
marrow or just solid. Overweight people have a tendency toward bone damage due to sustained compressions in bone joints and
tendons.
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Key Points
Elasticity is a measure of the deformation of an object when a force is applied. Objects that are very elastic like rubber have
high elasticity and stretch easily.
Stress is force over area.
Strain is change in length over original length.
Most objects behave elastically for small strains and return to their original shape after being bent.
If the strain on an object is greater than the elastic limit of the object, it will permanently deform or eventually fracture.
Fracture strength is a measure of the force needed to break an object.
Key Items
deformation: A transformation; change of shape.
strain: The amount by which a material deforms under stress or force, given as a ratio of the deformation to the initial
dimension of the material and typically symbolized by ε is termed the engineering strain. The true strain is defined as the
natural logarithm of the ratio of the final dimension to the initial dimension.
elastic: Capable of stretching; particularly, capable of stretching so as to return to an original shape or size when force is
released.
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8.7: The Center of Gravity
learning objectives
Describe how the center of mass of an oddly shaped object is found
Center of Gravity
When people think of objects, they think of them as singular particles of matter. In fact, every object is made up of millions of
particles, all of which behave differently when moved. When people observe a stick being thrown in the air, it seems as though the
entire object is moving at the same trajectory and velocity, but each particle is being subjected to a different motion in space and
acceleration, depending on its place. The different parts of the body have different motions. shows the motion of a stick in the air: it
seems to rotate around a single point. Three-dimensional bodies have a property called the center of mass, or center of gravity. This
center of mass’s main characteristic is that it appears to carry the whole mass of the body.
Center of Gravity: Although the center of mass is in the midpoint of the stick, all of the particles are moving as well.
The center of mass does not actually carry all the mass, despite appearances. Given a hollow sphere, the center is the center of
mass, even though it does not actually have anything in it. As seen in, it looks as if the external forces of gravity appear to be
working only on the center of mass, but each particle is being pushed or pulled by gravity. The center of mass is much easier to use
when discussing bodies, because no one has to analyze each individual particle.
Mathematical Expression: The mathematical relation of center of gravity is read as: ‘the position of the center of mass and
weighted average of the position of the particles. ‘
Specifically: ‘the total mass x the position of the center of mass= ∑ the mass of the individual particle x the position of the particle.
‘ The center of mass is a geometric point in three-dimensional volume. When using the definition above, it yields the following
equation for center of mass:
Σmi ri
rCOM = (8.7.1)
M
where r is the reference axis x, y, or z; m is individual mass;r is the individual position; and M M is the total mass.
i
When taking the center of mass of an oddly shaped object, it is helpful to break it down into smaller sections whose mass and
properties are easier to evaluate, and then add the products of the individual masses and positions and divide by the total mass.
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Center of Mass: This child’s toy uses the principles of ‘center of mass’ to stay balanced on a finger.
Key Points
The center of mass ‘s main characteristic is that it appears to carry the whole mass of the body.
The total mass x the position of the center of mass= ∑ mass of the individual particle x the position of the particle.
The center of mass is a geometric point in three-dimensional volume. By using the definition above, the following equation for
mi ri
center of mass can be derived: r =Σ
COM .
M
Key Terms
center of mass: The center of mass (COM) is the unique point at the center of a distribution of mass in space that has the
property that the weighted position vectors relative to this point sum to zero.
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8.8: Torque and Angular Acceleration
learning objectives
Express the relationship between the torque and the angular acceleration in a form of equation
Torque and angular acceleration are related by the following formula where is the objects moment of inertia and αα is the angular
acceleration.
Torque, Angular Acceleration, and the Role of the Church in the French Revolution: Why do things change their angular
velocity? Soon, you’ll know.
Just like Newton’s Second Law, which is force is equal to the mass times the acceleration, torque obeys a similar law. If you
replace torque with force and rotational inertia with mass and angular acceleration with linear acceleration, you get Newton’s
Second Law back out. In fact, this equation is Newton’s second law applied to a system of particles in rotation about a given axis. It
makes no assumptions about constant rotational velocity.
The net torque about an axis of rotation is equal to the product of the rotational inertia about that axis and the angular acceleration,
as shown in Figure 1.
8.8.1 https://phys.libretexts.org/@go/page/15505
Figure 1: Relationship between force (F), torque (τ), momentum (p), and angular momentum (L) vectors in a rotating system
Similar to Newton’s Second Law, angular motion also obeys Newton’s First Law. If no outside forces act on an object, an object in
motion remains in motion and an object at rest remains at rest. With rotating objects, we can say that unless an outside torque is
applied, a rotating object will stay rotating and an object at rest will not begin rotating.
If a turntable were spinning counter clockwise (when viewed from the top), and you applied your fingers to opposite sides the
turntable would begin to slow its spinning. From a translational viewpoint, at least, there would be no net force applied to the
turntable. The force that points to one side would be cancelled by the force that points to the other. The forces of the two fingers
would cancel. Therefore, the turntable would be in translational equilibrium. Despite that, the rotational velocity would be
decreased meaning that the acceleration would no longer be zero. From this we might conclude that just because a rotating object is
in translational equilibrium, it is not necessarily in rotational equilibrium.
Key Points
When a torque is applied to an object it begins to rotate with an acceleration inversely proportional to its moment of inertia.
This relation can be thought of as Newton’s Second Law for rotation. The moment of inertia is the rotational mass and the
torque is rotational force.
Angular motion obeys Newton’s First Law. If no outside forces act on an object, an object in motion remains in motion and an
object at rest remains at rest.
Key Terms
angular acceleration: The rate of change of angular velocity, often represented by α.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
rotational inertia: The tendency of a rotating object to remain rotating unless a torque is applied to it.
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CHAPTER OVERVIEW
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1
9.1: Quantities of Rotational Kinematics
learning objectives
Assess the relationship between radians the the revolution of a CD
When objects rotate about some axis—for example, when the CD (compact disc) rotates about its center—each point in the object
follows a circular arc. Consider a line from the center of the CD to its edge. Each pit used to record sound along this line moves
through the same angle in the same amount of time. The rotation angle is the amount of rotation, and is analogous to linear
distance. We define the rotation angle Δθ to be the ratio of the arc length to the radius of curvature:
Δθ =
Δs
r
(illustrated in ).
Rotation Angle: All points on a CD travel in circular arcs. The pits along a line from the center to the edge all move through the
same angle Δ in a time Δt.
In mathematics, the angle of rotation (or angular position ) is a measurement of the amount (i.e., the angle) that a figure is rotated
about a fixed point (often the center of a circle, as shown in ).
Angle θ and Arc Length s: The radius of a circle is rotated through an angle Δ. The arc length Δs is described on the
circumference.
The arc length Δs is the distance traveled along a circular path. r is the radius of curvature of the circular path. We know that for
one complete revolution, the arc length is the circumference of a circle of radius r. The circumference of a circle is 2πr. Thus, for
one complete revolution the rotation angle is:
(2πr)
Δθ = = 2π. (9.1.1)
r
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This result is the basis for defining the units used to measure rotation angles to be radians (rad), defined so that:
2π rad = 1 revolution. (9.1.2)
If Δθ = 2π rad, then the CD has made one complete revolution, and every point on the CD is back at its original position. Because
there are 360º in a circle or one revolution, the relationship between radians and degrees is thus 2π rad = 360º, so that:
360 º
1 rad = = 57.3 .º (9.1.3)
2π
Δt
.
learning objectives
Examine how fast an object is rotating based on angular velocity
To examine how fast an object is rotating, we define angular velocity ω as the rate of change of an angle. In symbols, this is
Δθ
ω = , (9.1.4)
Δt
where an angular rotation Δ takes place in a time Δt. The greater the rotation angle in a given amount of time, the greater the
angular velocity. The units for angular velocity are radians per second (rad/s).
Angular velocity ω is analogous to linear velocity v. To find the precise relationship between angular and linear velocity, we again
consider a pit on the rotating CD. This pit moves an arc length Δs in a time Δt, and so it has a linear velocity v = . Δs
Δt
(Δs) (r⋅Δθ)
From Δθ = r
we see that Δs = r ⋅ Δθ . Substituting this into the expression for v gives v = (Δt)
= r(
Δθ
Δt
) = rω.
r
.
The first relationship states that the linear velocity v is proportional to the distance from the center of rotation, thus it is largest for a
point on the rim (largest r), as you might expect. We can also call this linear speed v of a point on the rim the tangential speed. The
second relationship can be illustrated by considering the tire of a moving car, as shown in the picture below. Note that the speed of
the point at the center of the tire is the same as the speed v of the car. The faster the car moves, the faster the tire spins—large v
means a large ω, because v=rω. Similarly, a larger-radius tire rotating at the same angular velocity (ω) will produce a greater linear
speed (v) for the car.
Angular Velocity: A car moving at a velocity v to the right has a tire rotating with an angular velocity ω. The speed of the tread of
the tire relative to the axle is v, the same as if the car were jacked up. Thus the car moves forward at linear velocity v=rω, where r is
the tire radius. A larger angular velocity for the tire means a greater velocity for the car.
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Angular Acceleration, Alpha
Angular acceleration is the rate of change of angular velocity, expressed mathematically as α = Δω
Δt
.
learning objectives
Explain the relationship between angular acceleration and angular velocity
Angular acceleration is the rate of change of angular velocity. In SI units, it is measured in radians per second squared (rad/s2), and
is usually denoted by the Greek letter alpha (α ).
Consider the following situations in which angular velocity is not constant: when a skater pulls in her arms, when a child starts up a
merry-go-round from rest, or when a computer’s hard disk slows to a halt when switched off. In all these cases, there is an angular
acceleration in which ωω changes. The faster the change occurs, the greater the angular acceleration. Angular acceleration is
defined as the rate of change of angular velocity. In equation form, angular acceleration is expressed as follows:
Δω
α = (9.1.5)
Δt
where Δω is the change in angular velocity and Δt is the change in time. The units of angular acceleration are (rad/s)/s, or rad/s2.
If ωω increases, then αα is positive. If ωω decreases, then αα is negative.
It is useful to know how linear and angular acceleration are related. In circular motion, there is acceleration that is tangent to the
circle at the point of interest (as seen in the diagram below). This acceleration is called tangential acceleration, at.
Tangential acceleration: In circular motion, acceleration can occur as the magnitude of the velocity changes: a is tangent to the
motion. This acceleration is called tangential acceleration.
Tangential acceleration refers to changes in the magnitude of velocity but not its direction. In circular motion, centripetal
acceleration, ac, refers to changes in the direction of the velocity but not its magnitude. An object undergoing circular motion
experiences centripetal acceleration (as seen in the diagram below.) Thus, at and ac are perpendicular and independent of one
another. Tangential acceleration at is directly related to the angular acceleration and is linked to an increase or decrease in the
velocity (but not its direction).
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Centripetal Acceleration: Centripetal acceleration occurs as the direction of velocity changes; it is perpendicular to the circular
motion. Centripetal and tangential acceleration are thus perpendicular to each other.
Key Points
The arc length Δs is the distance traveled along a circular path. r is the radius of curvature of the circular path.
The rotation angle is the amount of rotation and is analogous to linear distance. We define the rotation angle Δθ to be the ratio
of the arc length to the radius of curvature: Δθ = Δs
r
.
For one complete revolution the rotation angle is 2π.
The greater the rotation angle in a given amount of time, the greater the angular velocity.
Angular velocity ω is analogous to linear velocity v.
We can write the relationship between linear velocity and angular velocity in two different ways: v = rω or ω = . v
The faster the change in angular velocity occurs, the greater the angular acceleration.
In circular motion, linear acceleration is tangent to the circle at the point of interest, and is called tangential acceleration.
In circular motion, centripetal acceleration refers to changes in the direction of the velocity but not its magnitude. An object
undergoing circular motion experiences centripetal acceleration.
Key Terms
Angular position: The angle in radians (degrees, revolutions) through which a point or line has been rotated in a specified
sense about a specified axis.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
angular acceleration: The rate of change of angular velocity, often represented by α.
tangential acceleration: The acceleration in a direction tangent to the circle at the point of interest in circular motion.
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9.10: Conservation of Energy
learning objectives
Conclude the interchangeability of force and radius with torque and angle of rotation in determining force
In this atom we will discuss work and energy associated with rotational motion. shows a worker using an electric grindstone
propelled by a motor. Sparks are flying, and noise and vibration are created as layers of steel are pared from the pole. The stone
continues to turn even after the motor is turned off, but it is eventually brought to a stop by friction. Clearly, the motor had to work
to get the stone spinning. This work went into heat, light, sound, vibration, and considerable rotational kinetic energy.
Grindstone: The motor works in spinning the grindstone, giving it rotational kinetic energy. That energy is then converted to heat,
light, sound, and vibration. (Credit: U.S. Navy photo by Mass Communication Specialist Seaman Zachary David Bell. )
Work must be done to rotate objects such as grindstones or merry-go-rounds. The simplest rotational situation is one in which the
net force is exerted perpendicular to the radius of a disc and remains perpendicular as the disc starts to rotate. The force is parallel
to the displacement, and so the net work (W) done is the product of the force (F) and the radius (r) of the disc (this is otherwise
known as torque(τ)) times the angle (θ) of rotation:
Work and energy in rotational motion are completely analogous to work and energy in translational motion and completely
transferrable. Just as in translational motion (where kinetic energy equals 1/2mv2 where m is mass and v is velocity ), energy is
conserved in rotational motion. Kinetic energy (K.E.) in rotational motion is related to moment of rotational inertia (I) and angular
velocity (ω):
1 2
KE = Iω . (9.10.2)
2
The final rotational kinetic energy equals the work done by the torque:
1 2
W = τθ = Iω = KE. (9.10.3)
2
This confirms that the work done went into rotational kinetic energy. To return to the grindstone example, work was done to give
the grindstone rotational energy, and work is done by friction so that it loses kinetic energy. However, the energy is never
destroyed; it merely changes form from rotation of the grindstone to heat when friction is applied.
Key Points
Rotating objects have rotational kinetic energy.
Rotational kinetic energy can change form if work is done on the object.
Energy is never destroyed, if rotational energy is gained or lost, something must have done work on it to change the form of the
energy.
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Key Terms
work: A measure of energy expended in moving an object; most commonly, force times displacement. No work is done if the
object does not move.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
rotational inertia: The tendency of a rotating object to remain rotating unless a torque is applied to it.
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9.2: Angular Acceleration
learning objectives
Relate angle of rotation, angular velocity, and angular acceleration to their equivalents in linear kinematics
Simply by using our intuition, we can begin to see the interrelatedness of rotational quantities like θ (angle of rotation), ω(angular
velocity) and α (angular acceleration). For example, if a motorcycle wheel has a large angular acceleration for a fairly long time, it
ends up spinning rapidly and rotating through many revolutions. The wheel’s rotational motion is analogous to the fact that the
motorcycle’s large translational acceleration produces a large final velocity, and the distance traveled will also be large.
Kinematic Equations
Kinematics is the description of motion. We have already studied kinematic equations governing linear motion under constant
acceleration:
v = v0 + at (9.2.1)
1 2
x = v0 t + at (9.2.2)
2
2 2
v =v + 2ax (9.2.3)
0
Similarly, the kinematics of rotational motion describes the relationships among rotation angle, angular velocity, angular
acceleration, and time. Let us start by finding an equation relating ω, α, and t. To determine this equation, we use the corresponding
equation for linear motion:
v = v0 + at. (9.2.4)
As in linear kinematics where we assumed a is constant, here we assume that angular acceleration α is a constant, and can use the
relation: a = rα Where r – radius of curve.Similarly, we have the following relationships between linear and angular values:
v = rω (9.2.5)
x = rθ (9.2.6)
By using the relationships a = rα, v = rω, and x = rθ , we derive all the other kinematic equations for rotational motion under
constant acceleration:
ω = ω0 + αt (9.2.7)
1
2
θ = ω0 t + αt (9.2.8)
2
2 2
ω =ω + 2αθ (9.2.9)
0
The equations given above can be used to solve any rotational or translational kinematics problem in which a and α are constant.
shows the relationship between some of the quantities discussed in this atom.
Linear and Angular: This figure shows uniform circular motion and some of its defined quantities.
9.2.1 https://phys.libretexts.org/@go/page/14488
Key Points
The kinematic equations for rotational and/or linear motion given here can be used to solve any rotational or translational
kinematics problem in which a and α are constant.
By using the relationships between velocity and angular velocity, distance and angle of rotation, and acceleration and angular
acceleration, rotational kinematic equations can be derived from their linear motion counterparts.
To derive rotational equations from the linear counterparts, we used the relationships a = rα, v = rω, and x = rθ .
Key Terms
kinematics: The branch of mechanics concerned with objects in motion, but not with the forces involved.
angular: Relating to an angle or angles; having an angle or angles; forming an angle or corner; sharp-cornered; pointed; as in,
an angular figure.
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9.3: Rotational Kinematics
learning objectives
Distinguish the two different motions in which rolling without slipping is broken down
Rolling without slipping generally occurs when an object rolls without skidding. To relate this to a real world phenomenon, we can
consider the example of a wheel rolling on a flat, horizontal surface.
Connecting Linear and Rotational Motion! Rolling without Slipping!: How fast does the axle of a bike wheel move? How fast
does the BOTTOM of a wheel move?
Rolling without slipping can be better understood by breaking it down into two different motions: 1) Motion of the center of mass,
with linear velocity v (translational motion); and 2) rotational motion around its center, with angular velocity w.
9.3.1 https://phys.libretexts.org/@go/page/14489
When an object is rolling on a plane without slipping, the point of contact of the object with the plane does not move. If we
imagine a wheel moving forward by rolling on a plane at speed v, it must also be rotating about its axis at an angular speed ωω
since it is rolling.
The object’s angular velocity ωω is directly proportional to its velocity, because as we know, the faster we are driving in a car, the
faster the wheels have to turn. So, to determine the exact relationship between linear velocity and angular velocity, we can imagine
the scenario in which the wheel travels a distance of x while rotating through an angle θ.
x = Rθ (9.3.1)
Since xx and θθ depend on time, we can take the derivative of both sides to obtain:
dx dθ
=R (9.3.2)
dt dt
where dx
dt
is equal to the linear velocity v , and dθ
dt
is equal to the angular velocity ωω. We can then simplify this equation to:
v = ωR (9.3.3)
Key Points
Rolling without slipping can be better understood by breaking it down into translational motion and rotational motion.
When an object is rolling on a plane without slipping, the point of contact of the object with the plane does not move.
A rolling object’s velocity v is directly related to its angular velocity ω, and is mathematically expressed as v = ωR , where R is
the object’s radius and v is its linear velocity.
Key Terms
angular velocity: A vector quantity describing the motion of an object in circular motion; its magnitude is equal to the angular
speed (ωω) of the particle, and the direction is perpendicular to the plane of its circular motion.
linear velocity: A vector quantity that denotes the rate of change of position with respect to time of the object’s center of mass.
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9.4: Dynamics
learning objectives
Explain the relationship between the force, mass, radius, and angular acceleration
If you have ever spun a bike wheel or pushed a merry-go-round, you have experienced the force needed to change angular velocity.
Our intuition is reliable in predicting many of the factors that are involved. For example, we know that a door opens slowly if we
push too close to its hinges. Furthermore, we know that the more massive the door, the more slowly it opens. The first example
implies that the farther the force is applied from the pivot, the greater the angular acceleration; another implication is that angular
acceleration is inversely proportional to mass. These relationships should seem very similar to the familiar relationships among
force, mass, and acceleration embodied in Newton’s second law of motion. There are, in fact, precise rotational analogs to both
force and mass.
Rotational inertia, as illustrated in, is the resistance of objects to changes in their rotation. In other words, a rotating object will stay
rotating and a non-rotating object will stay non-rotating unless acted on by a torque. This should remind you of Newton’s First
Law.
Rotational Inertia: Force is required to spin the bike wheel. The greater the force, the greater the angular acceleration produced.
The more massive the wheel, the smaller the angular acceleration. If you push on a spoke closer to the axle, the angular
acceleration will be smaller.
To develop the precise relationship among force, mass, radius, and angular acceleration, consider what happens if we exert a force
F on a point mass m that is at a distance r from a pivot point. Because the force is perpendicular to r, an acceleration a = is m
F
obtained in the direction of F. We can rearrange this equation such that F=ma and then look for ways to relate this expression to
expressions for rotational quantities. We note that a=rα, and we substitute this expression into F=ma, yielding:
F = mrα. (9.4.1)
Recall that torque is the turning effectiveness of a force. In this case, because F is perpendicular to r, torque is simply τ = Fr . So, if
we multiply both sides of the equation above by r, we get torque on the left-hand side. That is,
2
rF = mr α (9.4.2)
, or
2
τ = mr α. (9.4.3)
This equation is the rotational analog of Newton’s second law (F = ma ), where torque is analogous to force, angular acceleration
is analogous to translational acceleration, and mr2 is analogous to mass (or inertia). The quantity mr2 is called the rotational inertia
or moment of inertia of a point mass m a distance r from the center of rotation.
Different shapes of objects have different rotational inertia which depend on the distribution of their mass.
9.4.1 https://phys.libretexts.org/@go/page/14490
Key Points
The farther the force is applied from the pivot, the greater the angular acceleration.
Angular acceleration is inversely proportional to mass.
The equation τ = mr α. is the rotational analog of Newton’s second law (F = ma ), where torque is analogous to force,
2
angular acceleration is analogous to translational acceleration, and mr is analogous to mass (or inertia ).
2
Key Terms
rotational inertia: The tendency of a rotating object to remain rotating unless a torque is applied to it.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
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9.5: Rotational Kinetic Energy
learning objectives
Express the rotational kinetic energy as a function of the angular velocity and the moment of inertia, and relate it to the total
kinetic energy
Rotational kinetic energy is the kinetic energy due to the rotation of an object and is part of its total kinetic energy. Looking at
rotational energy separately around an object’s axis of rotation yields the following dependence on the object’s moment of inertia:
Kinetic Energy of Rotation: Things that roll without slipping have some fraction of their energy as translational kinetic and the
remainder as rotational kinetic. The ratio depends on the moment of inertia of the object that’s rolling.
1
2
Erotational = Iω , (9.5.1)
2
where ω is the angular velocity and I is the moment of inertia around the axis of rotation.
The mechanical work applied during rotation is the torque (τ ) times the rotation angle (θ) : W = τθ .
The instantaneous power of an angularly accelerating body is the torque times the angular velocity: P = τω .
Note the close relationship between the result for rotational energy and the energy held by linear (or translational) motion:
1
2
Etranslational = mv . (9.5.2)
2
9.5.1 https://phys.libretexts.org/@go/page/14491
In the rotating system, the moment of inertia takes the role of the mass and the angular velocity takes the role of the linear velocity.
As an example, let us calculate the rotational kinetic energy of the Earth (animated in Figure 1 ). As the Earth has a period of about
23.93 hours, it has an angular velocity of 7.29×10−5 rad/s. The Earth has a moment of inertia, I = 8.04×1037 kg·m2. Therefore, it
has a rotational kinetic energy of 2.138×1029 J.
The Rotating Earth: The earth’s rotation is a prominent example of rotational kinetic energy.
This can be partially tapped using tidal power. Additional friction of the two global tidal waves creates energy in a physical
manner, infinitesimally slowing down Earth’s angular velocity. Due to conservation of angular momentum this process transfers
angular momentum to the Moon’s orbital motion, increasing its distance from Earth and its orbital period.
Moment of Inertia
The moment of inertia is a property of a mass that measures its resistance to rotational acceleration about one or more axes.
learning objectives
Identify a property of a mass described by the moment of inertia
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AP Physics C - Moment of Inertia
Moment of Inertia: A brief introduction to moment of inertia (rotational inertia) for calculus-based physics students.
The moment of inertia I of an object can be defined as the sum of mr for all the point masses of which it is composed, where m is
2
the mass and r is the distance of the mass from the center of mass. It can be expressed mathematically as: I = ∑ mr . Here, I is
2
Moment of inertia also depends on the axis about which you rotate an object. Objects will usually rotate about their center of mass,
but can be made to rotate about any axis. The moment of inertia in the case of rotation about a different axis other than the center of
mass is given by the parallel axis theorem. The theorem states that the moment of inertia for an object rotated about a different axis
parallel to the axis passing through the center of mass is I + mr where r is now the distance between the two axes and IcmIcmis
cm
2
the moment of inertia when rotated about the center of mass which you learned how to calculate in the previous paragraph.
(net τ)
A general relationship among the torque, moment of inertia, and angular acceleration is: net τ = Iα, or α = I
. Net τ is the
total torque from all forces relative to a chosen axis. Such torques are either positive or negative and add like ordinary numbers.
The relationship in τ = Iα is the rotational analog to Newton’s second law and is very applicable. This equation is actually valid
for any torque, applied to any object, and relative to any axis.
As can be expected, the larger the torque, the larger the angular acceleration. For example, the harder a child pushes on a merry-go-
round, the slower it accelerates for the same torque. The basic relationship between the moment of inertia and the angular
acceleration is that the larger the moment of inertia, the smaller the angular acceleration. The moment of inertia depends not only
9.5.3 https://phys.libretexts.org/@go/page/14491
on the mass of an object, but also on its distribution of mass relative to the axis around which it rotates. For example, it wil be
much easier to accelerate a merry-go-round full of children if they stand close to its axis than if they all stand at the outer edge.
Moment of Inertia on a Merry-Go-Round: A father pushes a playground merry-go-round at its edge and perpendicular to its
radius to achieve maximum torque.
Key Points
1
Rotational kinetic energy can be expressed as: E rotational =
2
Iω where ω is the angular velocity and I is the moment of inertia
2
around the axis of rotation.
The mechanical work applied during rotation is the torque times the rotation angle: (θ) : W = τθ .
The instantaneous power of an angularly accelerating body is the torque times the angular velocity: P = τω .
There is a close relationship between the result for rotational energy and the energy held by linear (or translational) motion.
Newton’s first law, which describes the inertia of a body in linear motion, can be extended to the inertia of a body rotating about
an axis using the moment of inertia.
An object that is rotating at constant angular velocity will remain rotating unless it is acted upon by an external torque.
The larger the torque, the larger the angular acceleration.
Key Items
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
inertia: The property of a body that resists any change to its uniform motion; equivalent to its mass.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
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9.6: Conservation of Angular Momentum
learning objectives
Evaluate the implications of net torque on conservation of energy
Let us consider some examples of momentum: the Earth continues to spin at the same rate it has for billions of years; a high-diver
who is “rotating” when jumping off the board does not need to make any physical effort to continue rotating, and indeed would be
unable to stop rotating before hitting the water. These examples have the hallmarks of a conservation law. Following are further
observations to consider:
1. A closed system is involved. Nothing is making an effort to twist the Earth or the high-diver. They are isolated from rotation
changing influences (hence the term “closed system”).
2. Something remains unchanged. There appears to be a numerical quantity for measuring rotational motion such that the total
amount of that quantity remains constant in a closed system.
3. Something can be transferred back and forth without changing the total amount. A diver rotates faster with arms and legs
pulled toward the chest from a fully stretched posture.
Angular Momentum
The conserved quantity we are investigating is called angular momentum. The symbol for angular momentum is the letter L. Just as
linear momentum is conserved when there is no net external forces, angular momentum is constant or conserved when the net
torque is zero. We can see this by considering Newton’s 2nd law for rotational motion:
→ →
→
τ =
dL
dt
, where τ is the torque. For the situation in which the net torque is zero, dL
dt
=0 .
If the change in angular momentum ΔL is zero, then the angular momentum is constant; therefore,
→
L = constant (when net τ=0).
This is an expression for the law of conservation of angular momentum.
Conservation of Angular Momentum: An ice skater is spinning on the tip of her skate with her arms extended. Her angular
momentum is conserved because the net torque on her is negligibly small. In the next image, her rate of spin increases greatly when
she pulls in her arms, decreasing her moment of inertia. The work she does to pull in her arms results in an increase in rotational
kinetic energy.
→
→ →
(Both F and r are small, and so τ = r × F is negligibly small. ) Consequently, she can spin for quite some time. She can also
increase her rate of spin by pulling in her arms and legs. When she does this, the rotational inertia decreases and the rotation rate
increases in order to keep the angular momentum L = Iω constant. (I : rotational inertia, ω: angular velocity)
Conservation of angular momentum is one of the key conservation laws in physics, along with the conservation laws for energy
and (linear) momentum. These laws are applicable even in microscopic domains where quantum mechanics governs; they exist due
9.6.1 https://phys.libretexts.org/@go/page/14492
to inherent symmetries present in nature.
Rotational Collisions
In a closed system, angular momentum is conserved in a similar fashion as linear momentum.
learning objectives
Evaluate the difference in equation variables in rotational versus angular momentum
During a collision of objects in a closed system, momentum is always conserved. This fact is readily seen in linear motion. When
an object of mass m and velocity v collides with another object of mass m2 and velocity v2, the net momentum after the collision,
mv1f + mv2f, is the same as the momentum before the collision, mv1i + mv2i.
What if an rotational component of motion is introduced? Is momentum still conserved ?
9.6.2 https://phys.libretexts.org/@go/page/14492
Bowling ball and pi: When a bowling ball collides with a pin, linear and angular momentum is conserved
Yes. For objects with a rotational component, there exists angular momentum. Angular momentum is defined, mathematically, as
L = Iω , or L = r × p . This equation is an analog to the definition of linear momentum as p = mv . Units for linear momentum are
kg⋅m/s while units for angular momentum are kg⋅m2/s. As we would expect, an object that has a large moment of inertia I, such as
Earth, has a very large angular momentum. An object that has a large angular velocity ω, such as a centrifuge, also has a rather
large angular momentum.
So rotating objects that collide in a closed system conserve not only linear momentum p in all directions, but also angular
momentum L in all directions.
For example, take the case of an archer who decides to shoot an arrow of mass m1 at a stationary cylinder of mass m2 and radius r,
lying on its side. If the archer releases the arrow with a velocity v1i and the arrow hits the cylinder at its radial edge, what’s the final
momentum ?
Arrow hitting cyclinde: The arrow hits the edge of the cylinder causing it to roll.
Initially, the cylinder is stationary, so it has no momentum linearly or radially. Once the arrow is released, it has a linear momentum
p=m 1i and an angular component relative to the cylinders rotating axis, L = rp = rm v . After the collision, the arrow sticks to
1 1i
the rolling cylinder and the system has a net angular momentum equal to the original angular momentum of the arrow before the
collision.
Key Points
When an object is spinning in a closed system and no external torques are applied to it, it will have no change in angular
momentum.
The conservation of angular momentum explains the angular acceleration of an ice skater as she brings her arms and legs close
to the vertical axis of rotation.
If the net torque is zero, then angular momentum is constant or conserved.
Angular momentum is defined, mathematically, as L = Iω , or L = r × p . Which is the moment of inertia times the angular
velocity, or the radius of the object crossed with the linear momentum.
In a closed system, angular momentum is conserved in all directions after a collision.
Since momentum is conserved, part of the momentum in a collision may become angular momentum as an object starts to spin
after a collision.
Key Items
quantum mechanics: The branch of physics that studies matter and energy at the level of atoms and other elementary particles;
it substitutes probabilistic mechanisms for classical Newtonian ones.
9.6.3 https://phys.libretexts.org/@go/page/14492
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
angular momentum: A vector quantity describing an object in circular motion; its magnitude is equal to the momentum of the
particle, and the direction is perpendicular to the plane of its circular motion.
momentum: (of a body in motion) the product of its mass and velocity.
rotation: The act of turning around a centre or an axis.
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9.7: Vector Nature of Rotational Kinematics
learning objectives
Identify the direction of a vector using the Right Hand Rule
Angular momentum and angular velocity have both magnitude and direction and, therefore, are vector quantities. The direction of
these quantities is inherently difficult to track—a point on a rotating wheel is constantly rotating and changing direction. The axis
of rotation of a rotating wheel is the only place that has a fixed direction. The direction of angular momentum and velocity can be
determined along this axis.
Imagine the axis of rotation as a pole through the center of a wheel. The pole protrudes on both sides of the wheel and, depending
on which side you’re looking at, the wheel is turning either clockwise or counterclockwise. This dependency on perspective makes
determining the angle of rotation slightly more difficult. As with all physical quantities, there is a standard for measurement that
makes these types of quantities consistent. For angular quantities, the direction of the vector is determined using the Right Hand
Rule, illustrated in.
The Right Hand Rule: Figure (a) shows a disk is rotating counterclockwise when viewed from above. Figure (b) shows the right-
hand rule. The direction of angular velocity ω size and angular momentum L are defined to be the direction in which the thumb of
your right hand points when you curl your fingers in the direction of the disk’s rotation as shown.
The right hand rule can be used to find the direction of both the angular momentum and the angular velocity. From a spinning disc,
for example, let’s again imagine a pole through the center of the disc, at the axis of rotation. Using the right hand rule, your right
hand would be grasping the pole so that your four fingers (index, middle, ring, and pinky) are following the direction of rotation.
That is, an imaginary arrow from your wrist to your fingertips points in the same direction as the disc is rotating. In addition, your
thumb is pointing straight out in the axis, perpendicular to your other fingers (or parallel to the ‘pole’ at the axis of rotation). Using
this right hand rule, the direction of angular velocity ω and angular momentum L are defined as the direction in which the thumb of
your right hand points when you curl your fingers in the direction of the disc’s rotation.
Gyroscopes
A gyroscope is a spinning wheel or disk in which the axle is free to assume any orientation.
learning objectives
Compare the concept of a rotating wheel with a gyroscope
A gyroscope is a device for measuring or maintaining orientation based on the principles of angular momentum. Mechanically, a
gyroscope is a spinning wheel or disk in which the axle is free to assume any orientation. Although this orientation does not remain
fixed, it changes in response to an external torque much less and in a different direction than it would without the large angular
momentum associated with the disk’s high rate of spin and moment of inertia. The device’s orientation remains nearly fixed,
regardless of the mounting platform’s motion, because mounting the device in a gimbal minimizes external torque.
9.7.1 https://phys.libretexts.org/@go/page/15513
This equation means that the direction of ΔL is the same as the direction of the torque that creates it, as illustrated in. This direction
can be determined using the right hand rule, which says that the fingers on your hand curl towards the direction of rotation or force
exerted, and your thumb points towards the direction of angular momentum, torque, and angular velocity.
Direction of Torque and Angular Momentum: In figure (a), the torque is perpendicular to the plane formed by r and F and is the
direction your right thumb would point to if you curled your fingers in the direction of F. Figure (b) shows that the direction of the
torque is the same as that of the angular momentum it produces.
Rotating wheel: Consider a bicycle wheel with handles attached to it, as in. With the wheel rotating as shown, its angular
momentum is to the woman’s left. Suppose the person holding the wheel tries to rotate it as in the figure. Her natural expectation is
that the wheel will rotate in the direction she pushes it, however, what happens is quite different. The forces exerted create a torque
that is horizontal toward the person, and this torque creates a change in angular momentum L in the same direction, perpendicular
to the original angular momentum L, thus changing the direction of L but not the magnitude of L. ΔL and L add, giving a new
angular momentum with direction that is inclined more toward the person than before. The axis of the wheel has thus moved
perpendicular to the forces exerted on it, instead of in the expected direction.
Gyroscopic Effect: In figure (a), a person holding the spinning bike wheel lifts it with her right hand and pushes down with her left
hand in an attempt to rotate the wheel. This action creates a torque directly toward her. This torque causes a change in angular
momentum ΔL in exactly the same direction. Figure (b) shows a vector diagram depicting how ΔL and L add, producing a new
angular momentum pointing more toward the person. The wheel moves toward the person, perpendicular to the forces she exerts on
it.
Gyroscope: This same logic explains the behavior of gyroscopes (see ). There are two forces acting on a spinning gyroscope. The
torque produced is perpendicular to the angular momentum, thus the direction of the angular momentum is changed, but not its
9.7.2 https://phys.libretexts.org/@go/page/15513
magnitude. The gyroscope precesses around a vertical axis, since the torque is always horizontal and perpendicular to L. If the
gyroscope is not spinning, it acquires angular momentum in the direction of the torque (L=ΔL), and it rotates around a horizontal
axis, falling over just as we would expect.
Gyroscopes: As seen in figure (a), the forces on a spinning gyroscope are its weight and the supporting force from the stand. These
forces create a horizontal torque on the gyroscope, which create a change in angular momentum ΔL that is also horizontal. In figure
(b), ΔL and L add to produce a new angular momentum with the same magnitude, but different direction, so that the gyroscope
precesses in the direction shown instead of falling over.
Applications
Gyroscopes serve as rotational sensors. For this reason, applications of gyroscopes include inertial navigation systems where
magnetic compasses would not work (as in the Hubble telescope) or would not be precise enough (as in ICBMs). Another
application is the stabilization of flying vehicles, such as radio-controlled helicopters or unmanned aerial vehicles.
Key Points
Angular velocity and angular momentum are vector quantities and have both magnitude and direction.
The direction of angular velocity and angular momentum are perpendicular to the plane of rotation.
Using the right hand rule, the direction of both angular velocity and angular momentum is defined as the direction in which the
thumb of your right hand points when you curl your fingers in the direction of rotation.
Torque is perpendicular to the plane formed by r and F and is the direction your right thumb would point if you curled the
fingers of your right hand in the direction of F.
The direction of the torque is thus the same as that of the angular momentum it produces.
The gyroscope precesses around a vertical axis, since the torque is always horizontal and perpendicular to L. If the gyroscope is
not spinning, it acquires angular momentum in the direction of the torque, and it rotates about a horizontal axis, falling over just
as we would expect.
Key Items
angular momentum: A vector quantity describing an object in circular motion; its magnitude is equal to the momentum of the
particle, and the direction is perpendicular to the plane of its circular motion.
right hand rule: Direction of angular velocity ω and angular momentum L in which the thumb of your right hand points when
you curl your fingers in the direction of rotation.
angular velocity: A vector quantity describing an object in circular motion; its magnitude is equal to the speed of the particle
and the direction is perpendicular to the plane of its circular motion.
gimbal: A device for suspending something, such as a ship’s compass, so that it will remain level when its support is tipped.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
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9.8: Problem Solving
learning objectives
Develop and apply a strong problem-solving strategy for rotational kinematics
Example 9.8.1:
Suppose a large freight train accelerates from rest, giving its 0.350 m radius wheels an angular acceleration of 0.250 rad/s2.
After the wheels have made 200 revolutions (assume no slippage): (a) How far has the train moved down the track? (b) What
are the final angular velocity of the wheels and the linear velocity of the train?
In part (a), we are asked to find x, and in (b) we are asked to find ω and v. We are given the number of revolutions θ, the radius
of the wheels r, and the angular acceleration α.
The distance x is very easily found from the relationship between distance and rotation angle: θ = x
r
.
Solving this equation for x yields x = rθ .
Before using this equation, we must convert the number of revolutions into radians, because we are dealing with a relationship
between linear and rotational quantities:
2π rad
θ = (200rev)( ) = 1257 rad. (9.8.1)
1 rev
Substitute the known values into x = rθ to find the distance the train moved down the track:
x = rθ = (0.350 m)(1257 rad) = 440 m. (9.8.2)
We cannot use any equation that incorporates t to find ω, because the equation would have at least two unknown values. The
equation ω = ω + 2αθ will work, because we know the values for all variables except ω. Taking the square root of this
2 2
0
One may find the linear velocity of the train, v, through its relationship to ω:
v = rω = (0.350 m)(25.1 rad/s) = 8.77 m/s (9.8.5)
9.8.1 https://phys.libretexts.org/@go/page/15514
Joy physics probs 07- Rotational motion
Rotational motion: Part of a series of videos on physics problem-solving. The problems are taken from “The Joy of Physics. ”
This one deals with angular motion. The viewer is urged to pause the video at the problem statement and work the problem before
watching the rest of the video.
Key Points
Examine the situation to determine that rotational kinematics (rotational motion ) is involved, and identify exactly what needs to
be determined.
Make a list of what is given or can be inferred from the problem as stated and solve the appropriate equations.
9.8.2 https://phys.libretexts.org/@go/page/15514
Substitute the known values along with their units into the appropriate equation, and obtain numerical solutions complete with
units.
Key Terms
kinematics: The branch of mechanics concerned with objects in motion, but not with the forces involved.
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9.9: Linear and Rotational Quantities
learning objectives
Derive uniform circular motion from linear equations
A Rotating Body: Each particle constituting the body executes a uniform circular motion about the fixed axis. For the description
of the motion, angular quantities are the better choice.
When we describe the uniform circular motion in terms of angular velocity, there is no contradiction. The velocity (i.e. angular
velocity) is indeed constant. This is the first advantage of describing uniform circular motion in terms of angular velocity.
Second advantage is that angular velocity conveys the physical sense of the rotation of the particle as against linear velocity, which
indicates translational motion. Alternatively, angular description emphasizes the distinction between two types of motion
(translational and rotational).
Because dr
dt
= 0 for a uniform circular motion, we get v = ωr . Similarly, we also get a = αr where a stands for linear acceleration,
while α refers to angular acceleration (In a more general case, the relationship between angular and linear quantities are given as
v = ω × r, a = α × r + ω × v
1 1
2 2
θ = ω0 t + ( )α t : x = v0 t + ( )a t (9.9.3)
2 2
2 2 2 2
ω =ω + 2αθ : v =v + 2ax (9.9.4)
0 0
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Mass, Momentum, Energy, and Newton’s Second Law
As we use mass, linear momentum, translational kinetic energy, and Newton’s 2nd law to describe linear motion, we can describe a
general rotational motion using corresponding scalar/vector/tensor quantities:
Mass/ Rotational inertia:
Linenar/angular momentum:
Force/ Torque:
Kinetic energy:
For example, just as we use the equation of motion F=maF=ma to describe a linear motion, we can use its counterpart
τ =
dL
dt
= r×F to describe an angular motion. The descriptions are equivalent, and the choice can be made purely for the
convenience of use.
Key Points
As we use mass, linear momentum, translational kinetic energy, and Newton’s 2nd law to describe linear motion, we can
describe a general rotational motion using corresponding scalar/vector/tensor quantities.
Angular and linear velocity have the following relationship: v = ω × r .
As we use the equation of motion F = ma to describe a linear motion, we can use its counterpart τ = = r × F , to describe
dL
dt
angular motion. The descriptions are equivalent, and the choice can be made purely for the convenience of use.
Key Items
uniform circular motion: Movement around a circular path with constant speed.
torque: A rotational or twisting effect of a force; (SI unit newton-meter or Nm; imperial unit foot-pound or ft-lb)
rotational inertia: The tendency of a rotating object to remain rotating unless a torque is applied to it.
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Index
A angular frequency B
aberration 15.3: Periodic Motion ballistics
24.3: Lenses Angular momentum 5.4: Types of Forces in Nature
25.2: Other Optical Instruments 29.3: Atomic Physics and Quantum Mechanics basal metabolic rate
absolute space 5.9: Angular vs. Linear Quantities
9.6: Conservation of Angular Momentum 6.6: Power
27.1: Introduction 9.7: Vector Nature of Rotational Kinematics base
absolute zero angular motion 17.6: Applications of Electrostatics
12.2: Temperature and Temperature Scales 8.1: Introduction battery
12.4: Ideal Gas Law
14.4: Entropy angular position 19.2: Electric Current
9.1: Quantities of Rotational Kinematics 20.1: Overview
14.5: The Third Law of Thermodynamics
acceleration angular velocity Bell's theorem
5.3: Velocity, Acceleration, and Force 28.1: History and Quantum Mechanical Quantities
1.2: Units
2.3: Acceleration 5.6: Kepler’s Laws Bernoulli’s equation
2.5: Free-Falling Objects 5.9: Angular vs. Linear Quantities 11.3: Bernoulli’s Equation
3.2: Vectors 8.1: Introduction
9.1: Quantities of Rotational Kinematics
Beta decay
4.3: Newton’s Laws
5.3: Velocity, Acceleration, and Force 9.10: Conservation of Energy 30.2: Radioactivity
achromatic 9.3: Rotational Kinematics bilateral symmetry
9.5: Rotational Kinetic Energy 3.3: Projectile Motion
25.2: Other Optical Instruments 9.7: Vector Nature of Rotational Kinematics
achromatic doublet binary data
annihilation 26.4: Applications of Wave Optics
24.3: Lenses 27.3: Relativistic Quantities
action potential black body
antimatter 13.4: Methods of Heat Transfer
19.6: Electricity in the World 27.3: Relativistic Quantities 28.1: History and Quantum Mechanical Quantities
adhesion antinode 29.1: Overview
11.3: Bernoulli’s Equation 16.1: Introduction 29.2: The Early Atom
Adiabatic 16.5: Further Topics Black body radiation
12.4: Ideal Gas Law aperture 28.1: History and Quantum Mechanical Quantities
adiabatic index 25.1: The Human Eye blackbody
13.2: Specific Heat 25.2: Other Optical Instruments 13.5: Global Warming
aether application Boltzmann constant
27.2: Consequences of Special Relativity 1.1: The Basics of Physics 12.5: Kinetic Theory
afocal system approximation boundary condition
24.3: Lenses 1.3: Significant Figures and Order of Magnitude 15.6: Wave Behavior and Interaction
alpha particle Archimedes' Principle Boyle's Law
29.2: The Early Atom 10.3: Archimedes’ Principle 14.2: The First Law of Thermodynamics
30.2: Radioactivity artificial satellite brachytherapy
alternating current 5.6: Kepler’s Laws 30.4: Applications of Nuclear Physics
20.5: RC Circuits asteroid breakdown
Alveoli 14.4: Entropy 18.4: Capacitors and Dielectrics
10.2: Density and Pressure 5.1: Introduction to UCM and Gravitation
brine
AM Radio Waves astronomical unit 12.2: Temperature and Temperature Scales
23.1: The Electromagnetic Spectrum 5.6: Kepler’s Laws
Brownian motion
ampere atom 12.1: Introduction
19.2: Electric Current 12.1: Introduction 12.4: Ideal Gas Law
21.5: Magnetic Fields, Magnetic Forces, and atomic number buoyancy
Conductors 29.5: Multielectron Atoms 10.2: Density and Pressure
Ampere's law atomic orbital buoyant force
21.1: Magnetism and Magnetic Fields 28.2: Applications of Quantum Mechanics 10.3: Archimedes’ Principle
amplifier 29.5: Multielectron Atoms
22.3: Applications of Induction and EM Waves atomic spectra C
amplitude 30.1: The Nucleus
15.3: Periodic Motion atrium calorie
15.7: Waves on Strings 14.2: The First Law of Thermodynamics
11.4: Other Applications
16.2: Sound Intensity and Level
Avogadro’s number calorimeter
16.4: Interactions with Sound Waves 13.2: Specific Heat
Aneroid Barometer 12.4: Ideal Gas Law
29.2: The Early Atom calorimetry
10.2: Density and Pressure 13.1: Introduction
axis
angular Capacitance
3.2: Vectors
9.2: Angular Acceleration 18.4: Capacitors and Dielectrics
axon
angular acceleration
18.2: Equipotential Surfaces and Lines
8.8: Torque and Angular Acceleration 18.3: Point Charge
9.1: Quantities of Rotational Kinematics
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capacitor component Coulomb force
17.4: The Electric Field Revisited 3.2: Vectors 21.3: Magnetic Force on a Moving Electric Charge
18.4: Capacitors and Dielectrics compression 6.5: Potential Energy and Conservation of Energy
20.3: Kirchhoff’s Rules Coulomb’s Law
15.5: Waves
20.5: RC Circuits
compressive stress 10.6: Deformation of Solids
Carnot cycle 17.4: The Electric Field Revisited
8.5: Applications of Statics
14.4: Entropy Coulomb's Law
cathode computed tomography
17.3: Coulomb’s Law
30.4: Applications of Nuclear Physics
29.2: The Early Atom crest
CCD concave lens
15.5: Waves
24.3: Lenses
29.4: Applications of Atomic Physics Critical temperature
cell membrane concave mirror
19.3: Resistance and Resistors
24.4: Mirrors
18.2: Equipotential Surfaces and Lines critically damped
18.3: Point Charge concentration
15.4: Damped and Driven Oscillations
center of mass 12.9: Diffusion
cryptography
5.5: Newton’s Law of Universal Gravitation conductive medium
28.1: History and Quantum Mechanical Quantities
7.5: Center of Mass 19.1: Overview
8.3: Stability conductor crystallography
8.7: The Center of Gravity 25.2: Other Optical Instruments
17.1: Overview
centrifugal force 17.4: The Electric Field Revisited
29.2: The Early Atom
5.4: Types of Forces in Nature 18.2: Equipotential Surfaces and Lines Curie temperature
centripetal 18.3: Point Charge 21.2: Magnets
18.4: Capacitors and Dielectrics Current
29.2: The Early Atom
23.3: Applications of EM Waves
5.1: Introduction to UCM and Gravitation 19.2: Electric Current
5.2: Non-Uniform Circular Motion conservation 20.3: Kirchhoff’s Rules
5.3: Velocity, Acceleration, and Force 6.5: Potential Energy and Conservation of Energy 21.5: Magnetic Fields, Magnetic Forces, and
centripetal acceleration 7.1: Introduction Conductors
15.3: Periodic Motion Conservation of Energy Cyclotron
CFC 7.3: Collisions 21.4: Motion of a Charged Particle in a Magnetic
conservation of momentum principle Field
14.3: The Second Law of Thermodynamics
characteristic time constant 7.3: Collisions cyclotron frequency
conservative force 21.4: Motion of a Charged Particle in a Magnetic
22.2: AC Circuits
Field
chemical reaction 6.5: Potential Energy and Conservation of Energy
12.1: Introduction constant velocity
D
Chromatic aberration 3.1: Motion in Two Dimensions
25.2: Other Optical Instruments constructive interference damping
16.5: Further Topics
circuit 15.6: Wave Behavior and Interaction
18.4: Capacitors and Dielectrics
16.5: Further Topics de Broglie wavelength
25.2: Other Optical Instruments 29.4: Applications of Atomic Physics
Circular Motion 26.2: Diffraction
5.3: Velocity, Acceleration, and Force
Decay
continuity 30.2: Radioactivity
cladding 10.5: Fluids in Motion
24.2: Reflection, Refraction, and Dispersion
decibels (units)
contraction 16.2: Sound Intensity and Level
classical electrodynamics 11.4: Other Applications
23.2: Electromagnetic Waves and their Properties
deformation
contrast 15.2: Hooke’s Law
Classical mechanics 25.2: Other Optical Instruments 8.6: Elasticity, Stress, Strain, and Fracture
27.3: Relativistic Quantities 26.4: Applications of Wave Optics
degeneracy
classical physics control rod 14.5: The Third Law of Thermodynamics
16.3: Doppler Effect and Sonic Booms 30.4: Applications of Nuclear Physics
degree of freedom
closed system conversion 12.1: Introduction
7.1: Introduction 1.2: Units 7.3: Collisions
cochlea convex demagnetization
16.2: Sound Intensity and Level 25.2: Other Optical Instruments 14.5: The Third Law of Thermodynamics
coherence convex lens destructive interference
28.2: Applications of Quantum Mechanics 24.3: Lenses 15.6: Wave Behavior and Interaction
coherent convex mirror 16.5: Further Topics
16.4: Interactions with Sound Waves 24.4: Mirrors 25.2: Other Optical Instruments
26.1: Superposition and Interference 26.2: Diffraction
coordinate
cohesion diagram
3.2: Vectors
11.3: Bernoulli’s Equation 2.4: Problem-Solving for Basic Kinematics
Coordinate axes
combination circuit dialysis
3.2: Vectors
20.2: Resistors in Series and Parallel 11.1: Overview
corona discharge
combustion diamagnetism
17.6: Applications of Electrostatics
13.2: Specific Heat 21.6: Applications of Magnetism
coulomb
complex numbers diastole
17.1: Overview
19.5: Alternating Currents 18.1: Overview 11.2: Flow in Tubes
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dielectric drift velocity electromotive force (emf)
17.1: Overview 19.2: Electric Current 20.3: Kirchhoff’s Rules
17.2: Shelding and Charging Through Induction 21.5: Magnetic Fields, Magnetic Forces, and 22.1: Magnetic Flux, Induction, and Faraday’s Law
17.4: The Electric Field Revisited Conductors electron affinity
18.4: Capacitors and Dielectrics dynamic contrast ratio 28.2: Applications of Quantum Mechanics
24.2: Reflection, Refraction, and Dispersion
25.1: The Human Eye electron gun
differential dynamics 18.5: Applications
12.8: Thermal Stresses
1.4: Solving Physics Problems electron shell
differential equation 4.7: Further Applications of Newton’s Laws
21.6: Applications of Magnetism
20.5: RC Circuits dynamo 29.5: Multielectron Atoms
23.2: Electromagnetic Waves and their Properties
21.1: Magnetism and Magnetic Fields electron volt
diffraction
18.1: Overview
25.2: Other Optical Instruments
26.2: Diffraction
E electrostatic force
28.1: History and Quantum Mechanical Quantities eardrum 17.3: Coulomb’s Law
29.2: The Early Atom 16.2: Sound Intensity and Level element
diffusion eccentricity 29.5: Multielectron Atoms
11.4: Other Applications 5.6: Kepler’s Laws elementary charge
12.9: Diffusion
19.6: Electricity in the World
elastic 21.5: Magnetic Fields, Magnetic Forces, and
7.2: Conservation of Momentum Conductors
dimension 8.6: Elasticity, Stress, Strain, and Fracture elementary particle
1.4: Solving Physics Problems
10.6: Deformation of Solids
Elastic Collision 28.1: History and Quantum Mechanical Quantities
7.3: Collisions 7.1: Introduction emissivity
7.3: Collisions
diode 13.4: Methods of Heat Transfer
30.2: Radioactivity
elastic potential energy 23.1: The Electromagnetic Spectrum
30.4: Applications of Nuclear Physics 15.3: Periodic Motion Energy
diopter elasticity 15.5: Waves
25.2: Other Optical Instruments 15.2: Hooke’s Law 6.1: Introduction
dipole 16.1: Introduction Enthalpy
24.2: Reflection, Refraction, and Dispersion
electric charge 10.1: Introduction
17.1: Overview 13.2: Specific Heat
dipole moment 17.5: Electric Flux and Gauss’s Law entropy
17.1: Overview 19.1: Overview
18.1: Overview 12.2: Temperature and Temperature Scales
20.3: Kirchhoff’s Rules 14.3: The Second Law of Thermodynamics
21.6: Applications of Magnetism 21.6: Applications of Magnetism 14.4: Entropy
direct current electric current epistemological
20.1: Overview 19.2: Electric Current
20.5: RC Circuits 28.1: History and Quantum Mechanical Quantities
electric displacement field epithelium
direction of propagation 17.5: Electric Flux and Gauss’s Law
15.5: Waves 16.4: Interactions with Sound Waves
Electric field equation
discharge 17.1: Overview
17.1: Overview 4.5: Problem-Solving
17.5: Electric Flux and Gauss’s Law
17.2: Shelding and Charging Through Induction 18.1: Overview equilibrium
disorder 21.5: Magnetic Fields, Magnetic Forces, and 10.2: Density and Pressure
14.4: Entropy Conductors 12.6: Phase Changes
22.4: Magnetic Fields and Maxwell Revisited 12.9: Diffusion
Dispersion 13.6: Phase Equilbrium
26.3: Further Topics
electric potential
15.4: Damped and Driven Oscillations
18.1: Overview
displacement 17.4: The Electric Field Revisited
16.4: Interactions with Sound Waves
electrical circuit 4.7: Further Applications of Newton’s Laws
20.1: Overview 8.2: Conditions for Equilibrium
2.1: Basics of Kinematics
3.2: Vectors electrical resistance equipotential
dissipative force 19.1: Overview 18.1: Overview
18.2: Equipotential Surfaces and Lines
15.3: Periodic Motion electrocardiogram 18.3: Point Charge
6.5: Potential Energy and Conservation of Energy 19.6: Electricity in the World
distortion escape velocity
electromagnet 5.8: Energy Conservation
25.2: Other Optical Instruments 21.1: Magnetism and Magnetic Fields
diurnal Euclidean
electromagnetic force 10.6: Deformation of Solids
5.4: Types of Forces in Nature 29.2: The Early Atom
Doppler effect evaporation
electromagnetic radiation 13.6: Phase Equilbrium
16.3: Doppler Effect and Sonic Booms 26.3: Further Topics
23.2: Electromagnetic Waves and their Properties 30.2: Radioactivity
exponent
Doppler shift 6.8: Further Topics 1.3: Significant Figures and Order of Magnitude
29.2: The Early Atom Electromagnetic waves external beam therapy
dosimeter 23.2: Electromagnetic Waves and their Properties 30.4: Applications of Nuclear Physics
30.4: Applications of Nuclear Physics electromotive force
dot product 20.1: Overview
6.2: Work Done by a Constant Force 20.2: Resistors in Series and Parallel
22.4: Magnetic Fields and Maxwell Revisited
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F frequency ground state
Faraday constant 15.3: Periodic Motion 28.2: Applications of Quantum Mechanics
15.5: Waves group velocity
12.4: Ideal Gas Law 16.1: Introduction
Faraday shield 15.5: Waves
16.3: Doppler Effect and Sonic Booms
17.6: Applications of Electrostatics 16.4: Interactions with Sound Waves gyroradius
22.3: Applications of Induction and EM Waves 21.4: Motion of a Charged Particle in a Magnetic
Faraday's law of induction 23.2: Electromagnetic Waves and their Properties Field
20.1: Overview
22.1: Magnetic Flux, Induction, and Faraday’s Law
friction
ferromagnet 4.3: Newton’s Laws H
4.7: Further Applications of Newton’s Laws
22.1: Magnetic Flux, Induction, and Faraday’s Law 6.5: Potential Energy and Conservation of Energy
harmonic oscillator
ferromagnetic 7.3: Collisions 28.1: History and Quantum Mechanical Quantities
21.1: Magnetism and Magnetic Fields frigorific mixture heat
21.6: Applications of Magnetism 12.2: Temperature and Temperature Scales 14.1: Introduction
fibrillation fundamental thermodynamic relation 14.2: The First Law of Thermodynamics
19.6: Electricity in the World 13.2: Specific Heat
Heat capacity
fictitious force fusion 13.2: Specific Heat
5.4: Types of Forces in Nature 30.3: Quantum Tunneling and Conservation Laws
heat engine
field of view 30.4: Applications of Nuclear Physics 14.3: The Second Law of Thermodynamics
6.8: Further Topics 14.4: Entropy
25.1: The Human Eye
first law of thermodynamics heat of reaction
12.4: Ideal Gas Law G 13.2: Specific Heat
14.3: The Second Law of Thermodynamics Galilean transformation heat pump
6.8: Further Topics 14.4: Entropy
27.3: Relativistic Quantities
fission galvanometer heat transfer
30.3: Quantum Tunneling and Conservation Laws 13.1: Introduction
20.4: Voltmeters and Ammeters
30.4: Applications of Nuclear Physics helical
22.1: Magnetic Flux, Induction, and Faraday’s Law
6.8: Further Topics
fluid gamma (γ) rays 26.4: Applications of Wave Optics
23.1: The Electromagnetic Spectrum helical motion
4.7: Further Applications of Newton’s Laws
29.2: The Early Atom 21.4: Motion of a Charged Particle in a Magnetic
fluidity 30.2: Radioactivity Field
10.1: Introduction 30.4: Applications of Nuclear Physics
hertz
flux Gamma decay 16.1: Introduction
22.1: Magnetic Flux, Induction, and Faraday’s Law 23.1: The Electromagnetic Spectrum
Hohmann transfer orbit
23.2: Electromagnetic Waves and their Properties gamma ray 5.6: Kepler’s Laws
flux density 29.3: Atomic Physics and Quantum Mechanics
Hooke's Law
18.3: Point Charge gas constant 6.5: Potential Energy and Conservation of Energy
flux intensity 12.4: Ideal Gas Law
humidity
18.2: Equipotential Surfaces and Lines Gauge Pressure 12.6: Phase Changes
FM 10.2: Density and Pressure
Huygen's principle
23.1: The Electromagnetic Spectrum general relativity 16.5: Further Topics
focal point 27.4: Implications of Special Relativity
hydraulic press
24.3: Lenses geometric optics 10.2: Density and Pressure
Force 24.1: Overview
hydrogen bond
15.4: Damped and Driven Oscillations geothermal 12.3: Thermal Expansion
4.1: Introduction
14.4: Entropy
4.2: Force and Mass
6.3: Work Done by a Variable Force gimbal I
7.3: Collisions 9.7: Vector Nature of Rotational Kinematics
8.2: Conditions for Equilibrium ideal energy
gluon
fossil fuels 12.4: Ideal Gas Law
30.1: The Nucleus
6.7: CASE STUDY: World Energy Use ideal fluid
gradient
Fourier analysis 11.3: Bernoulli’s Equation
5.4: Types of Forces in Nature
15.1: Introduction Ideal gas
grating
frame of reference 10.2: Density and Pressure
28.1: History and Quantum Mechanical Quantities 12.1: Introduction
2.1: Basics of Kinematics Gravitational acceleration 12.2: Temperature and Temperature Scales
22.1: Magnetic Flux, Induction, and Faraday’s Law 12.4: Ideal Gas Law
4.4: Other Examples of Forces
free electron laser Gravitational Force 12.5: Kinetic Theory
28.2: Applications of Quantum Mechanics 14.2: The First Law of Thermodynamics
5.5: Newton’s Law of Universal Gravitation
free fall image distance
gravity 24.3: Lenses
2.5: Free-Falling Objects
17.1: Overview
3.3: Projectile Motion
impedance
5.7: Gravitational Potential Energy 20.5: RC Circuits
22.2: AC Circuits
greenhouse effect
13.5: Global Warming
impulse
7.1: Introduction
greenhouse gas
13.5: Global Warming
incandescence
26.4: Applications of Wave Optics
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incident ray isentropic lens
26.3: Further Topics 14.5: The Third Law of Thermodynamics 25.2: Other Optical Instruments
incline isolated system 26.1: Superposition and Interference
4.7: Further Applications of Newton’s Laws 13.1: Introduction Lenz's law
incompressible 6.5: Potential Energy and Conservation of Energy 22.2: AC Circuits
10.2: Density and Pressure isotope leverage
10.5: Fluids in Motion 30.4: Applications of Nuclear Physics 8.5: Applications of Statics
11.3: Bernoulli’s Equation isotropic line element
index of refraction 12.3: Thermal Expansion 27.4: Implications of Special Relativity
24.2: Reflection, Refraction, and Dispersion linear momentum
26.3: Further Topics
J 7.1: Introduction
induction linear thermal expansion coefficient
22.1: Magnetic Flux, Induction, and Faraday’s Law
joint
8.5: Applications of Statics 12.3: Thermal Expansion
5.5: Newton’s Law of Universal Gravitation
inductor linear velocity
17.2: Shelding and Charging Through Induction K 9.3: Rotational Kinematics
22.1: Magnetic Flux, Induction, and Faraday’s Law kelvin Longitudinal
22.2: AC Circuits 15.5: Waves
16.1: Introduction
22.4: Magnetic Fields and Maxwell Revisited Lorentz factor
inelastic Kelvin scale
12.2: Temperature and Temperature Scales 27.2: Consequences of Special Relativity
7.1: Introduction 27.3: Relativistic Quantities
7.2: Conservation of Momentum kilocalorie
Lorentz force
inelastic scattering 13.1: Introduction
17.3: Coulomb’s Law
29.2: The Early Atom Kinematic 22.1: Magnetic Flux, Induction, and Faraday’s Law
inertia 2.3: Acceleration Lorentz transformations
2.4: Problem-Solving for Basic Kinematics
1.2: Units 27.1: Introduction
3.1: Motion in Two Dimensions
4.3: Newton’s Laws 27.3: Relativistic Quantities
9.2: Angular Acceleration
9.5: Rotational Kinetic Energy 27.4: Implications of Special Relativity
9.8: Problem Solving
inertial frame kinematics Lorenz Invariance
5.4: Types of Forces in Nature 27.1: Introduction
2.1: Basics of Kinematics
instantaneous Kinetic Energy luminance
2.2: Speed and Velocity 25.1: The Human Eye
10.2: Density and Pressure
insulator 15.2: Hooke’s Law luminiferous aether
17.1: Overview 18.1: Overview 27.1: Introduction
interfere 4.7: Further Applications of Newton’s Laws
5.8: Energy Conservation
16.4: Interactions with Sound Waves 6.8: Further Topics M
interference 7.3: Collisions Mach number
15.6: Wave Behavior and Interaction kinetic theory of gases 16.3: Doppler Effect and Sonic Booms
26.2: Diffraction
26.3: Further Topics
12.1: Introduction machines
12.5: Kinetic Theory
26.4: Applications of Wave Optics 8.5: Applications of Statics
28.2: Applications of Quantum Mechanics magnetic domains
interferometer L 21.2: Magnets
26.1: Superposition and Interference laminar Magnetic Field
27.3: Relativistic Quantities 11.4: Other Applications 21.3: Magnetic Force on a Moving Electric Charge
intermolecular laminar flow 21.5: Magnetic Fields, Magnetic Forces, and
10.4: Cohesion and Adhesion 11.1: Overview Conductors
12.6: Phase Changes 21.6: Applications of Magnetism
land 22.1: Magnetic Flux, Induction, and Faraday’s Law
internal energy 26.4: Applications of Wave Optics 22.4: Magnetic Fields and Maxwell Revisited
12.4: Ideal Gas Law laser
13.1: Introduction
magnetic field lines
14.1: Introduction 28.2: Applications of Quantum Mechanics 21.1: Magnetism and Magnetic Fields
14.2: The First Law of Thermodynamics latent heat of fusion magnetic flux
14.3: The Second Law of Thermodynamics 13.3: Phase Change and Latent Heat 22.1: Magnetic Flux, Induction, and Faraday’s Law
International Systems of Units latent heat of vaporization magnetic mirror
12.4: Ideal Gas Law 13.3: Phase Change and Latent Heat 21.4: Motion of a Charged Particle in a Magnetic
inverse law Field
5.5: Newton’s Law of Universal Gravitation 1.1: The Basics of Physics magnetron
inviscid law of conservation of energy 21.4: Motion of a Charged Particle in a Magnetic
11.3: Bernoulli’s Equation 14.2: The First Law of Thermodynamics Field
ionization energy LCD Magnification
28.2: Applications of Quantum Mechanics 26.4: Applications of Wave Optics 24.3: Lenses
ionizing radiation length magnitude
23.1: The Electromagnetic Spectrum 1.2: Units 3.2: Vectors
30.4: Applications of Nuclear Physics
length contraction mass
iridescence 27.4: Implications of Special Relativity 21.6: Applications of Magnetism
26.2: Diffraction 4.2: Force and Mass
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mass distribution N optical window
15.3: Periodic Motion nanostructure 23.1: The Electromagnetic Spectrum
Mass Spectrometer 25.2: Other Optical Instruments order of magnitude
21.4: Motion of a Charged Particle in a Magnetic natural convection 1.3: Significant Figures and Order of Magnitude
Field origin
13.4: Methods of Heat Transfer
matter 3.2: Vectors
natural frequency
1.1: The Basics of Physics orthogonal
16.5: Further Topics
matter waves 17.4: The Electric Field Revisited
natural satellite
28.1: History and Quantum Mechanical Quantities 21.4: Motion of a Charged Particle in a Magnetic
29.2: The Early Atom 5.6: Kepler’s Laws Field
29.3: Atomic Physics and Quantum Mechanics nematic 26.1: Superposition and Interference
Maxwell’s Equations 26.4: Applications of Wave Optics oscillate
22.1: Magnetic Flux, Induction, and Faraday’s Law nerve impulse 15.2: Hooke’s Law
27.1: Introduction 16.4: Interactions with Sound Waves 15.3: Periodic Motion
28.1: History and Quantum Mechanical Quantities 15.7: Waves on Strings
29.2: The Early Atom
net force 26.3: Further Topics
4.3: Newton’s Laws
mean motion oscillator
5.6: Kepler’s Laws
neuron 15.4: Damped and Driven Oscillations
18.2: Equipotential Surfaces and Lines
mechanical advantage 18.3: Point Charge
osmosis
8.5: Applications of Statics 11.1: Overview
neurotransmitter
mechanical equivalent of heat 16.4: Interactions with Sound Waves
over damping
13.1: Introduction 15.4: Damped and Driven Oscillations
Newton’s third law of motion
media 7.4: Rocket Propulsion
oxidation
16.1: Introduction 14.2: The First Law of Thermodynamics
Newton's Law of Gravitation
medium 5.7: Gravitational Potential Energy
ozone
15.5: Waves 23.1: The Electromagnetic Spectrum
Newtonian mechanics
metabolism 12.5: Kinetic Theory
14.2: The First Law of Thermodynamics
NMR P
meteorology 30.4: Applications of Nuclear Physics parallel
5.4: Types of Forces in Nature 20.2: Resistors in Series and Parallel
noble gas
metric 12.5: Kinetic Theory parallel equivalent resistance
27.4: Implications of Special Relativity 19.3: Resistance and Resistors
node
microscopy 16.1: Introduction Paramagnetism
26.4: Applications of Wave Optics 16.5: Further Topics 21.6: Applications of Magnetism
microstate Nonconservative Forces particle accelerator
14.5: The Third Law of Thermodynamics 6.5: Potential Energy and Conservation of Energy 18.1: Overview
Minkowski space normal 28.2: Applications of Quantum Mechanics
27.4: Implications of Special Relativity 4.4: Other Examples of Forces pendentive
model normal force 8.5: Applications of Statics
1.1: The Basics of Physics 5.1: Introduction to UCM and Gravitation Pendulums
mole nuclear reactor 6.8: Further Topics
12.4: Ideal Gas Law 30.4: Applications of Nuclear Physics perihelion
moment of inertia nucleons 5.6: Kepler’s Laws
12.5: Kinetic Theory 30.3: Quantum Tunneling and Conservation Laws period
8.4: Solving Statics Problems 15.3: Periodic Motion
nucleosynthesis
Momentum 30.4: Applications of Nuclear Physics
16.1: Introduction
4.3: Newton’s Laws
nucleus periodic table
7.1: Introduction 29.5: Multielectron Atoms
7.3: Collisions 17.1: Overview
9.6: Conservation of Angular Momentum 29.2: The Early Atom permanent magnet
30.1: The Nucleus 21.1: Magnetism and Magnetic Fields
monochromatic
26.1: Superposition and Interference
nuclide Permeability
26.2: Diffraction 30.1: The Nucleus 22.1: Magnetic Flux, Induction, and Faraday’s Law
28.2: Applications of Quantum Mechanics null measurements permittivity
motion 20.4: Voltmeters and Ammeters 18.4: Capacitors and Dielectrics
2.4: Problem-Solving for Basic Kinematics perpendicular
Motional Emf O 4.4: Other Examples of Forces
4.7: Further Applications of Newton’s Laws
22.1: Magnetic Flux, Induction, and Faraday’s Law Objective measurement
MRI 16.1: Introduction
phase
23.2: Electromagnetic Waves and their Properties
30.4: Applications of Nuclear Physics ohm
muscles 19.2: Electric Current
phase velocity
15.5: Waves
8.5: Applications of Statics Ohm’s law
mutual inductance 22.2: AC Circuits
phasor
22.2: AC Circuits
22.2: AC Circuits ohmic
myocardium 19.3: Resistance and Resistors
Phasors
19.6: Electricity in the World 19.5: Alternating Currents
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phon Potential Energy rarefaction
16.2: Sound Intensity and Level 18.1: Overview 15.5: Waves
phosphor 5.7: Gravitational Potential Energy raster
5.8: Energy Conservation
18.5: Applications 18.5: Applications
6.8: Further Topics
photoconductivity potentiometer ray tracing
17.6: Applications of Electrostatics 24.3: Lenses
20.4: Voltmeters and Ammeters
Photoelectric effect power Rayleigh criterion
15.6: Wave Behavior and Interaction 28.1: History and Quantum Mechanical Quantities
15.5: Waves
28.1: History and Quantum Mechanical Quantities Reactance
6.6: Power
29.2: The Early Atom
photoelectron prefix 22.2: AC Circuits
1.2: Units reflection
28.1: History and Quantum Mechanical Quantities
photon pressure 24.1: Overview
24.2: Reflection, Refraction, and Dispersion
10.4: Cohesion and Adhesion
23.2: Electromagnetic Waves and their Properties 26.3: Further Topics
28.2: Applications of Quantum Mechanics Probability density function refraction
29.2: The Early Atom 28.1: History and Quantum Mechanical Quantities
24.1: Overview
physical pendulum propulsion 24.2: Reflection, Refraction, and Dispersion
15.3: Periodic Motion 5.8: Energy Conservation 25.2: Other Optical Instruments
pit pupil 26.3: Further Topics
26.4: Applications of Wave Optics 25.1: The Human Eye Refractive index
Planck constant Purkinje fibers 15.6: Wave Behavior and Interaction
23.2: Electromagnetic Waves and their Properties
29.1: Overview 19.6: Electricity in the World
Planck’s constant relative
3.4: Multiple Velocities
28.1: History and Quantum Mechanical Quantities Q
plane relativistic quantum theory
Quantization 28.1: History and Quantum Mechanical Quantities
10.6: Deformation of Solids 29.2: The Early Atom
Plane waves 29.3: Atomic Physics and Quantum Mechanics
relativity of simultaneity
27.4: Implications of Special Relativity
15.5: Waves quantum mechanics
26.1: Superposition and Interference 9.6: Conservation of Angular Momentum
renewable forms of energy
planet 6.7: CASE STUDY: World Energy Use
Quarks
5.1: Introduction to UCM and Gravitation 30.1: The Nucleus
reorientate
Plasma 3.3: Projectile Motion
10.1: Introduction R resisitivity
12.6: Phase Changes 17.1: Overview
radar 19.3: Resistance and Resistors
plumb line
23.1: The Electromagnetic Spectrum resistance
7.5: Center of Mass
radial 20.2: Resistors in Series and Parallel
point mass
18.1: Overview Resistor
5.5: Newton’s Law of Universal Gravitation 5.2: Non-Uniform Circular Motion
point particle 20.3: Kirchhoff’s Rules
Radians 20.5: RC Circuits
7.5: Center of Mass 5.3: Velocity, Acceleration, and Force
Poiseuille's Law resolution
radiation 26.2: Diffraction
10.2: Density and Pressure 1.2: Units
polarity resonance
radiative transfer 15.6: Wave Behavior and Interaction
17.6: Applications of Electrostatics 13.5: Global Warming 22.2: AC Circuits
polarizability radio waves rest mass
26.3: Further Topics 23.1: The Electromagnetic Spectrum 27.3: Relativistic Quantities
polarization 23.3: Applications of EM Waves restoring force
26.3: Further Topics Radioactive decay 15.3: Periodic Motion
polarizer 30.1: The Nucleus 15.6: Wave Behavior and Interaction
24.2: Reflection, Refraction, and Dispersion 30.2: Radioactivity resultant
30.4: Applications of Nuclear Physics
position 4.6: Vector Nature of Forces
radioactive tracer
2.3: Acceleration reversible
30.4: Applications of Nuclear Physics
positive feedback 14.2: The First Law of Thermodynamics
13.4: Methods of Heat Transfer
radiograph 14.4: Entropy
23.1: The Electromagnetic Spectrum Reynolds number
positron
27.3: Relativistic Quantities
radiography 11.2: Flow in Tubes
29.3: Atomic Physics and Quantum Mechanics 25.2: Other Optical Instruments right hand rule
30.2: Radioactivity radioisotope 21.3: Magnetic Force on a Moving Electric Charge
30.4: Applications of Nuclear Physics 30.2: Radioactivity 9.7: Vector Nature of Rotational Kinematics
potential radiometric dating rigid
12.3: Thermal Expansion 30.2: Radioactivity 4.7: Further Applications of Newton’s Laws
6.5: Potential Energy and Conservation of Energy
radionuclide rigid body
potential difference
30.1: The Nucleus 7.5: Center of Mass
18.1: Overview 30.2: Radioactivity
20.2: Resistors in Series and Parallel
radon
30.2: Radioactivity
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rms silver halide Static Equilibrium
12.5: Kinetic Theory 28.2: Applications of Quantum Mechanics 10.2: Density and Pressure
20.5: RC Circuits simple circuit 17.4: The Electric Field Revisited
22.2: AC Circuits 18.2: Equipotential Surfaces and Lines
19.3: Resistance and Resistors
rms current 18.3: Point Charge
Simple Harmonic Motion 8.3: Stability
19.5: Alternating Currents
rms voltage
15.6: Wave Behavior and Interaction stimulated emission
simple harmonic oscillator 29.4: Applications of Atomic Physics
19.5: Alternating Currents
root mean square
15.3: Periodic Motion Stoke's theorem
simple pendulum 22.1: Magnetic Flux, Induction, and Faraday’s Law
19.5: Alternating Currents
rotation
15.3: Periodic Motion strain
sinoatrial node 4.7: Further Applications of Newton’s Laws
9.6: Conservation of Angular Momentum
19.6: Electricity in the World 8.6: Elasticity, Stress, Strain, and Fracture
Rotational Inertia sinusoidal streamlined
8.8: Torque and Angular Acceleration
15.3: Periodic Motion 11.4: Other Applications
9.10: Conservation of Energy
9.4: Dynamics sinusoidal steady state stress
9.9: Linear and Rotational Quantities 19.5: Alternating Currents 12.8: Thermal Stresses
Snell’s law of refraction 4.7: Further Applications of Newton’s Laws
S 15.6: Wave Behavior and Interaction string theory
24.2: Reflection, Refraction, and Dispersion 28.1: History and Quantum Mechanical Quantities
scalar
2.1: Basics of Kinematics solenoid stroboscopic
3.2: Vectors 21.2: Magnets 2.4: Problem-Solving for Basic Kinematics
scanning tunneling microscope (STM) 22.1: Magnetic Flux, Induction, and Faraday’s Law Subjective measurement
28.2: Applications of Quantum Mechanics sonic boom 16.1: Introduction
Schrodinger's Equation 16.3: Doppler Effect and Sonic Booms sublimation
28.1: History and Quantum Mechanical Quantities Special relativity 10.1: Introduction
10.6: Deformation of Solids 13.3: Phase Change and Latent Heat
scientific method
1.1: The Basics of Physics
22.1: Magnetic Flux, Induction, and Faraday’s Law Superconductivity
23.2: Electromagnetic Waves and their Properties
19.3: Resistance and Resistors
scientific notation 26.4: Applications of Wave Optics
1.3: Significant Figures and Order of Magnitude 27.1: Introduction superimpose
27.2: Consequences of Special Relativity 16.4: Interactions with Sound Waves
Scintillation 27.3: Relativistic Quantities 25.2: Other Optical Instruments
29.2: The Early Atom 27.4: Implications of Special Relativity 26.4: Applications of Wave Optics
scintillators 28.1: History and Quantum Mechanical Quantities superposition
30.2: Radioactivity specific heat 15.6: Wave Behavior and Interaction
second law of thermodynamics 12.4: Ideal Gas Law 16.4: Interactions with Sound Waves
14.3: The Second Law of Thermodynamics 13.2: Specific Heat 17.4: The Electric Field Revisited
selective absorber spectral color surface vertices
13.5: Global Warming 23.1: The Electromagnetic Spectrum 24.3: Lenses
semiclassical spectral radiance symmetrical
29.2: The Early Atom 29.1: Overview 3.3: Projectile Motion
semiclassical approach spectrum symmetry
28.2: Applications of Quantum Mechanics 29.2: The Early Atom 4.3: Newton’s Laws
semiconductor speed of light systole
19.3: Resistance and Resistors 27.2: Consequences of Special Relativity 11.2: Flow in Tubes
30.2: Radioactivity 27.3: Relativistic Quantities
series speed of propagation T
15.5: Waves
20.2: Resistors in Series and Parallel tangent
series equivalent resistance spherical aberration 4.7: Further Applications of Newton’s Laws
25.2: Other Optical Instruments
19.3: Resistance and Resistors tangential acceleration
shear stress spin 9.1: Quantities of Rotational Kinematics
21.6: Applications of Magnetism
10.1: Introduction
28.2: Applications of Quantum Mechanics
telecommunication
11.2: Flow in Tubes
23.3: Applications of EM Waves
shock hazard stable equilibrium
8.3: Stability
temperature coefficient of resistivity
19.5: Alternating Currents
19.3: Resistance and Resistors
19.6: Electricity in the World standard atmosphere
Shunt resistance 12.2: Temperature and Temperature Scales
tensile stress
8.5: Applications of Statics
20.4: Voltmeters and Ammeters standing wave
shutter speed 15.7: Waves on Strings
terminal velocity
29.2: The Early Atom 17.1: Overview
25.2: Other Optical Instruments
SI units static terminal voltage
1.4: Solving Physics Problems 20.2: Resistors in Series and Parallel
12.4: Ideal Gas Law
sidereal year 4.7: Further Applications of Newton’s Laws tesla
5.6: Kepler’s Laws
static contrast ratio 21.3: Magnetic Force on a Moving Electric Charge
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theory Transformer vector field
1.1: The Basics of Physics 22.1: Magnetic Flux, Induction, and Faraday’s Law 10.6: Deformation of Solids
thermal agitation translation 17.4: The Electric Field Revisited
23.1: The Electromagnetic Spectrum 8.2: Conditions for Equilibrium velocity
thermal conductivity transmutation 2.2: Speed and Velocity
2.3: Acceleration
13.4: Methods of Heat Transfer 30.2: Radioactivity
3.2: Vectors
thermal energy transverse 4.2: Force and Mass
14.3: The Second Law of Thermodynamics 21.5: Magnetic Fields, Magnetic Forces, and 5.3: Velocity, Acceleration, and Force
thermal equilibrium Conductors ventricle
12.7: The Zeroth Law of Thermodynamics transverse wave 11.4: Other Applications
13.1: Introduction 15.5: Waves vessel
14.1: Introduction 15.7: Waves on Strings
11.2: Flow in Tubes
thermal hazard 16.5: Further Topics
trigonometry Vienna Standard Mean Ocean Water
19.5: Alternating Currents
12.4: Ideal Gas Law
19.6: Electricity in the World 1.4: Solving Physics Problems
thermal radiation triple point virtual image
24.4: Mirrors
23.1: The Electromagnetic Spectrum 12.2: Temperature and Temperature Scales
thermodynamic temperature trough viscosity
10.1: Introduction
14.1: Introduction 15.5: Waves
11.2: Flow in Tubes
Thermodynamics tunneling 11.3: Bernoulli’s Equation
12.2: Temperature and Temperature Scales 30.3: Quantum Tunneling and Conservation Laws 11.4: Other Applications
12.6: Phase Changes Turbine visible light
14.1: Introduction
22.1: Magnetic Flux, Induction, and Faraday’s Law 23.1: The Electromagnetic Spectrum
thermography turbulence voltage
23.1: The Electromagnetic Spectrum
11.4: Other Applications 17.1: Overview
Thick Lenses turbulent 19.2: Electric Current
24.3: Lenses
11.4: Other Applications
thin lens turbulent flow W
24.3: Lenses warfare
11.1: Overview
thin lens equation 30.4: Applications of Nuclear Physics
24.3: Lenses
U watt
Thomson scattering 6.6: Power
29.2: The Early Atom
Ultrasound
15.6: Wave Behavior and Interaction wave
Thoracic Cavity 15.5: Waves
10.2: Density and Pressure
underdamped
15.4: Damped and Driven Oscillations Wave equation
thrust 15.6: Wave Behavior and Interaction
4.3: Newton’s Laws
Uniform Circular Motion
15.3: Periodic Motion wave speed
time dilation 9.9: Linear and Rotational Quantities 15.5: Waves
27.2: Consequences of Special Relativity waveform
27.4: Implications of Special Relativity
uniform motion
4.3: Newton’s Laws 15.1: Introduction
tomography
unit vector wavefront
25.2: Other Optical Instruments
30.4: Applications of Nuclear Physics 17.3: Coulomb’s Law 16.5: Further Topics
torque 3.2: Vectors wavelength
10.2: Density and Pressure 15.5: Waves
18.1: Overview V 17.2: Shelding and Charging Through Induction
22.3: Applications of Induction and EM Waves
21.5: Magnetic Fields, Magnetic Forces, and valance shell 23.2: Electromagnetic Waves and their Properties
Conductors
29.5: Multielectron Atoms 26.1: Superposition and Interference
22.1: Magnetic Flux, Induction, and Faraday’s Law
6.4: Work-Energy Theorem vapor 26.3: Further Topics
7.5: Center of Mass 12.6: Phase Changes weight
8.2: Conditions for Equilibrium vaporization 5.5: Newton’s Law of Universal Gravitation
8.4: Solving Statics Problems
8.8: Torque and Angular Acceleration
13.6: Phase Equilbrium wettability
9.4: Dynamics vascular 11.3: Bernoulli’s Equation
9.5: Rotational Kinetic Energy 11.1: Overview Wheatstone bridge
9.6: Conservation of Angular Momentum vector 20.4: Voltmeters and Ammeters
9.7: Vector Nature of Rotational Kinematics
9.9: Linear and Rotational Quantities 17.4: The Electric Field Revisited work
18.1: Overview 15.5: Waves
torques 2.1: Basics of Kinematics 18.1: Overview
4.7: Further Applications of Newton’s Laws 3.2: Vectors 6.2: Work Done by a Constant Force
Torr 4.2: Force and Mass 6.3: Work Done by a Variable Force
4.6: Vector Nature of Forces 9.10: Conservation of Energy
10.2: Density and Pressure
5.9: Angular vs. Linear Quantities
tracer 8.1: Introduction
30.4: Applications of Nuclear Physics vector area X
trajectory 22.1: Magnetic Flux, Induction, and Faraday’s Law xerography
3.3: Projectile Motion 17.6: Applications of Electrostatics
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Z
zeroth law of thermodynamics
12.7: The Zeroth Law of Thermodynamics
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Glossary
Sample Word 1 | Sample Definition 1
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Detailed Licensing
Overview
Title: Book: Physics (Boundless)
Webpages: 207
All licenses found:
Undeclared: 99% (205 pages)
CC BY-SA 4.0: 1% (2 pages)
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