Christopher W. Kulp, Vasilis Pagonis - Classical Mechanics - A Computational Approach With Examples Using Mathematica and Python
Christopher W. Kulp, Vasilis Pagonis - Classical Mechanics - A Computational Approach With Examples Using Mathematica and Python
Classical Mechanics
i i
i i
i i
Classical Mechanics
A Computational Approach with
Examples Using Mathematica and
Python
Christopher W. Kulp
Vasilis Pagonis
i i
i i
i i
c 2021 Taylor & Francis Group, LLC
Reasonable efforts have been made to publish reliable data and information, but the author and publisher
cannot assume responsibility for the validity of all materials or the consequences of their use. The authors
and publishers have attempted to trace the copyright holders of all material reproduced in this publication
and apologize to copyright holders if permission to publish in this form has not been obtained. If any
copyright material has not been acknowledged please write and let us know so we may rectify in any future
reprint.
Except as permitted under U.S. Copyright Law, no part of this book may be reprinted, reproduced, trans-
mitted, or utilized in any form by any electronic, mechanical, or other means, now known or hereafter
invented, including photocopying, microfilming, and recording, or in any information storage or retrieval
system, without written permission from the publishers.
For permission to photocopy or use material electronically from this work, access www.copyright.com or
contact the Copyright Clearance Center, Inc. (CCC), 222 Rosewood Drive, Danvers, MA 01923, 978-750-
8400. For works that are not available on CCC please contact [email protected]
Trademark notice: Product or corporate names may be trademarks or registered trademarks and are used
only for identification and explanation without intent to infringe.
Typeset in LMRoman
by Nova Techset Private Limited, Bengaluru & Chennai, India
i i
i i
i i
Chris dedicates this book to his wife Gail, mother Linda, and his late father,
Chester. Without their support, this book would not have been possible.
Vasilis dedicates this book to his wife, Mary Jo Boylan, and to his students
at McDaniel College.
i i
i i
i i
Contents
Preface xiii
vii
i i
i i
i i
viii Contents
i i
i i
i i
Contents ix
i i
i i
i i
x Contents
i i
i i
i i
Contents xi
i i
i i
i i
xii Contents
Bibliography 445
Index 447
i i
i i
i i
Preface
xiii
i i
i i
i i
xiv Preface
[de Lange and Pierrus(2010), Taylor(2005), Morin(2008), Thornton and Marion(2004)] for
additional problems. These texts approach classical mechanics in different ways, and it is
beneficial to see multiple presentations of the same topic.
Technology has often disrupted the status quo on how things are done. When pocket
calculators were first introduced, people were afraid that it would ruin students’ ability to do
math. However a 2003 study found that “students’ operational skills and problem solving
skills improved when calculators were an integral part of instruction” [Ellington(2003)].
With further advances in technology, one can ask whether the use of computer algebra
systems (CAS) will damage a student’s ability to solve physics problems. A 2008 dissertation
[Tokpah(2008)] showed that “students using CAS tend to perform better than students
taught using non-CAS instruction” when learning mathematics. However, we believe that
the use of computers to assist, for example, in algebraic manipulations, will greatly benefit
physics students. Long and involved algebraic manipulation is rarely insightful in terms of
understanding the underlying physics. When a student is allowed to offload tedious algebraic
manipulations to a computer, this student can then focus on higher-level mathematics, such
as analyzing the problem using limits, in order to better develop a physical intuition of what
the equations are describing.
This book is for students who are taking a semester of classical mechanics, following
a full course in introductory physics. The prerequisites for this book are two semesters of
introductory physics and two semesters of calculus. Despite the title of the book, a semester
of computer science is not needed (but would be helpful) to start reading this book.
A NOTE TO INSTRUCTORS
One of our motivations for writing this book was to better prepare our students for the
large variety of careers that they pursue after graduating from a physics program. We have
found that after graduation, our students are increasingly taking on careers where compu-
tation is a critical element. Our students often take at least one computer science course.
i i
i i
i i
Preface xv
However, they were not always making the connection of how to apply the programming
skills they learned in computer science for the purpose of solving problems in physics. Part
of the reason is because traditional physics courses are still focused on closed-form solutions
and, when computation is used, it is generally focused on using computer algebra systems
to perform complicated integrals, with the occasional numerical solution of a differential
equation thrown in for good measure.
We believe that it is time for computation to be more closely integrated into the physics
curriculum. Doing so clearly demonstrates to students how to use computation to solve
problems, a skill that many of them will find critical whether they go to graduate school,
or enter into careers in government or industry after graduation. We should remember that
the majority of physics students (65% at the time of this writing [Supiano(2018)]) do not
go to graduate programs in physics or astronomy. Of the graduates who work in the private
sector, 77% work either in engineering, computer or information systems, or non-science
fields that regularly solve technical problems [Supiano(2018)]. Teaching students to use
computation to solve physics problems will provide them a transferable skill that they can
apply in their future careers, even if they never need to find a Lagrangian after taking their
classical mechanics course.
Of course including computation in a classical mechanics course (or any physics course)
comes at a cost, but we believe the payoff is well worth the cost. Every minute spent on
the instruction of programming and computation is one less minute that can be spent on
the instruction of topics like conservation laws, Hamiltonians, etc. In addition, by using
this textbook the student may be in a position where he/she needs to learn computer
programming, the use of a computer algebra system, and classical mechanics. This can be
a challenging position to be in as a student. However, today’s online resources can allow
students to more quickly pick up computing skills.
A quick online search will result in sites that provide online courses in Python pro-
gramming. However, even simple web searches such as, “How do I integrate an equation
in Mathematica?” can provide the student enough instruction for solving problems in both
Python and Mathematica. We find that the best method of learning how to use com-
putational tools is by actually solving problems with them. Physics students often learn
mathematics along with the physics topics, in order to solve physics problems. The same
can be true for programming.
At this point, you might be wondering about the trade-off in content? Both of the
authors teach one-semester–long classical mechanics courses and are aware of sacrifices that
need to be made when introducing new material in a well-established course like classical
mechanics. We did our best to include in this book all of the major topics traditionally
covered in a classical mechanics course. We believe we were successful in that endeavor.
This book can be used like any other classical mechanics book, such as the classic Classical
Dynamics of Particles and Systems by Thornton and Marion, and the chapters are in fact
structured like many other texts in the field.
An instructor can skip much of the discussion of computation and use this book like any
other. However, the inclusion of the computation allows the students to explore problems
that are difficult, or inaccessible without computation. It may be the case that in a one-
semester course, the instructor may need to be more selective in the physics topics covered.
We believe the exchange is worthwhile. The payoff is the ability to explore problems that
are not solvable in closed-form. Many real-world problems that students will encounter are
not solvable in closed form. Furthermore, the aforementioned transferable skill of using
computation to solve problems is an additional payoff that will serve students well for their
entire career.
i i
i i
i i
xvi Preface
ACKNOWLEDGMENTS
Finally, we want to thank several people who have helped make this book happen.We thank
Kirsten Barr, Rebecca Davies, and Shashi Kumar at CRC Press for all of their help in the
preparation of this book. We also extend our heart felt thanks to Mary Jo Boylan and
Maryam Esmat for their hard work and time they have put into this book. Mary Jo typed
the author’s hand written solutions for the instructor’s manual. She put in countless hours
preparing the solution manual to this book. We would have never made our deadline without
her! We very much appreciate her efforts. Maryam Esmat helped us proofread our book and
provided comments from a student’s perspective. Her detailed comments on each chapter
were invaluable to the development of this book. We believe the book is much better because
of Maryam’s efforts. We thank her for all of her time reading and rereading each chapter!
i i
i i
i i
CHAPTER 1
So, you have decided to—or are required to—learn the subject of classical mechanics. But,
what is classical mechanics? Is it fixing old cars? No, but a knowledge of classical mechanics
will help you understand how your car works! To understand the term, classical mechanics
let us first understand the term classical physics. Classical physics are the fields of physics
that don’t involve either quantum theory or the theory of relativity. Mechanics is the branch
of physics that deals with the actions of forces on an object that involve motion. So, classical
mechanics involves the study of forces on objects that are well-described without using
quantum theory and whose motion is nonrelativistic (i.e., cases where relativity is not
needed to correctly model the motion). In other words, classical mechanics is the physics of
your day-to-day world! An understanding of classical mechanics will help you understand
how to build roller coasters, merry-go-rounds, and airplanes. Our knowledge of classical
mechanics also allows us to make predictions on the motion of objects like comets and
punted footballs. The world you interact with on a daily basis is generally the world of
classical mechanics.
At first it might seem that classical mechanics is a dusty old subject that you need
to learn in order to get to the “interesting stuff” like quantum mechanics. While the field
of classical mechanics is one of the older subjects in physics, it is certainly not dusty! In
fact, there is a lot of intriguing current research done in classical systems. For example,
in Chapter 13 we will explore the field of nonlinear systems, which occur in in many of
the natural and social sciences, not only in physics, and they display interesting types of
behaviors including chaos.
1
i i
i i
i i
• Thermodynamics deals with relationships between all forms of energy; often the focus
is on heat and its relationship with other forms of energy.
• Electromagnetism deals with the interaction of electrically-charged particles using the
concepts of electric and magnetic fields.
• Statistical Mechanics deals with understanding how macroscopic properties, such as
temperature and pressure, emerge from a large number of particles that make up the
system.
• Relativity deals with the dependence of physical phenomenon on the relative motion
between the observer and the observed. Physicists often consider three cases of relativ-
ity: Galilean relativity (which falls under the category of classical physics), Einstein’s
theory of special relativity, and Einstein’s theory of general relativity. Special rela-
tivity focuses on the dependencies between inertial frames of reference while general
relativity, a generalization of special relativity, takes into account noninertial frames.
• Quantum Mechanics deals with the interactions between subatomic particles and
between subatomic particles and radiation.
Classical mechanics is one of the first topics learned by a physics student because it deals
with the more intuitive concepts of force and motion. It provides tools for not just describing
motion but also predicting motion. One of the central ideas of classical mechanics is that if
one knows the position and velocity of a particle at time, t, then one can find the particle’s
position and velocity at any point in the past or future (except for the case of chaotic
systems, which we will deal with later). The ability to make predictions is critical in science,
hence, classical mechanics provides a powerful set of tools, which are widely applicable to
many different types of systems.
As mentioned earlier, almost all real-world problems involve more than one branch
of physics, and classical mechanics is often an important element to those problems. For
example, one might want to know the motion of a charged particle in a magnetic field. The
first step would be to find the force on the particle from the magnetic field using formulas
from electromagnetism, then the position and velocity as a function of time can be found
using Newton’s second law (a formula from classical mechanics). Although this example is
simple, it illustrates the point that problems rarely involve only one branch of physics. Even
if one is working with a system in which quantum theory must be included, semiclassical
approaches, where part of the problem is treated classically, are sometimes quite useful.
It is common for classical mechanics to take a central role in many physics problems
because of the question it addresses.
Given the forces acting on an object, what is the resulting motion of that object?
Hopefully, you understand how classical mechanics fits in with the other branches and
its central importance to the larger field of physics. In the rest of this chapter, we will
lay down the important foundations for classical mechanics, starting with the next section
where we will explore the basic assumptions of classical mechanics.
i i
i i
i i
Classical mechanics addresses the problem of predicting an object’s motion given the
forces acting on the object. For now, we will consider the object to be a point particle. The
advantage of working with point particles is that they have no size and no internal dynam-
ics. In addition, point particles do not rotate nor do they deform, further simplifying their
dynamics. It turns out that treating objects as point particles can be a very good approx-
imation for describing translational motion. In later chapters, we will study the physics
of extended bodies and rotational dynamics, where we will no longer restrict ourselves to
working with point particles.
r = xî + y ĵ + z k̂ (1.2.1)
where î, ĵ, and k̂ are the Cartesian unit vectors along the x, y, and z axes, respectively. The
variables x, y, and z in (1.2.1) give the distance between the particle and the origin along
each axis and are called the components of the vector. Also note that in (1.2.1): r = r(t),
x = x(t), y = y(t), and z = z(t). In this book, vectors are denoted by bold font. In Chapter
3, we will discuss vector quantities in more detail.
The vector in (1.2.1) tells us that in order to get to the location of the particle, r, one
needs to move a distance x along the x-axis (as denoted by î), turn and move a distance y
along the y-axis (as denoted by ĵ), and then turn and move a distance z along the z-axis
i i
i i
i i
(as denoted by k̂). The notation used in this book is that a hat (the ˆ symbol) represents
a unit vector, a vector of length one. Hence, we can think of r as the sum of three vectors:
xî, y ĵ, and z k̂, each representing a displacement from the origin along one of the axes.
The displacement of the particle, ∆r = r − r0 , is the change of the particle’s position
from the position, r0 , to the position, r. The displacement is found by subtracting the two
vectors r0 and r, component by component,
There are other coordinate systems which we will explore later in this book. Each coor-
dinate system will have its own components and unit vectors, however, the basic idea of
position and displacement will be the same.
1.2.1.2 Velocity
The velocity v of a particle is the particle’s displacement (change of position) per unit
time and is measured in meters per second (m/s). The instantaneous velocity is found by
computing the time derivative of the position vector:
dr ∆r
v= = lim (1.2.4)
dt ∆t→0 ∆t
where ∆r = r(t + ∆t) − r(t). All of the vectors we will come across here are differentiable,
and the limits will exist. Furthermore, the derivative in (1.2.4) will behave like derivatives
you have encountered before. Therefore, we can use the standard derivative rules:
d dr1 dr2
(r1 + r2 ) = + (1.2.5)
dt dt dt
d dc dr
(cr) =r + c (1.2.6)
dt dt dt
where c is a scalar function of time. The rules in (1.2.5) and (1.2.6) allow us to distribute
the derivative to each vector component. In addition, if the unit vectors are constant, like
they are in Cartesian coordinates (but not in others!), then we can compute the velocity
by:
i i
i i
i i
dr
v= (1.2.7)
dt
d
= xî + y ĵ + z k̂ (1.2.8)
dt
dx dy dz
= î + ĵ + k̂ (1.2.9)
dt dt dt
or: v = vx î + vy ĵ + vz k̂ where vx = dx/dt, and so on. Note that there are only three terms
(and not six) in (1.2.9) because the Cartesian unit vectors are constant. Note that the speed
of a particle is the magnitude of its velocity vector. We will discuss how to calculate vector
magnitudes in Chapter 3.
Finally, we can further simplify the notation by using dots to denote differentiation with
respect to time (i.e., ẋ = dx/dt). Hence:
1.2.1.3 Acceleration
Acceleration, a, is the change of velocity per unit time and is measured in meters per second
squared (m/s2 ). The acceleration of an object is computed similarly to velocity:
dv
a= (1.2.11)
dt
d
= vx î + vy ĵ + vz k̂ (1.2.12)
dt
dvx dvy dvz
= î + ĵ + k̂ (1.2.13)
dt dt dt
=v̇x î + v̇y ĵ + v̇z k̂ (1.2.14)
which, like velocity can be rewritten as a = ax î + ay ĵ + az k̂. The acceleration is also the
second derivative of the position vector and using the dot notation we can write:
d2 r
a= (1.2.15)
dt2
= ẍî + ÿ ĵ + z̈ k̂ (1.2.16)
As we will see in future chapters, the acceleration is very important in classical mechan-
ics. Much of what is discussed in this book is about how to develop an equation for the
acceleration, which can then be integrated to get the motion of the system. All measure-
ments of position, velocity, and acceleration, are relative to the frame of reference (choice of
origin and axes) from which the measurements are made. We will return to the important
idea of reference frames in Section 1.4.
i i
i i
i i
1.2.2.1 Mass
The mass of an object is a measure of the object’s inertia, how strongly the object resists
acceleration. For example, a loaded shopping cart is more difficult to accelerate than an
empty one. Recall from your introductory physics course that mass is a scalar quantity
because in order to describe mass, only a magnitude is needed.
To compare the mass of objects, one could use a beam balance. In order to speak
quantitatively about mass, it is helpful to have a standard mass to which all other masses
could be compared. Other physical quantities, such as the meter and second, are based on
fundamental constants. As mentioned earlier, the meter is based on the speed of light, and
the second is based on the ground state hyperfine splitting frequency of Cesium-133 atoms,
respectively. The definition of the kilogram has recently changed in order to be based on a
fundamental constant.
Before 2019, the unit of mass, the kilogram, was defined to be the mass of a platinum-
iridium cylinder stored at the International Bureau of Weights and Measures outside of
Paris, France. Hence, the mass of an object could be found by placing the international
standard on one side of the beam balance and the object on the other. If the object has
a mass of one kilogram, then the balance will remain level, otherwise the balance will tilt
towards the more massive object.
The problem with using the old definition of the kilogram is that the mass of the standard
changed over time. Although the international standard and its copies (many countries have
several copies) were carefully stored, they were occasionally exposed to air and, therefore,
absorbed atmospheric contaminates. The problem here is that the copies used by various
countries gained mass at different rates than the international standard. Although these
differences might be in the micrograms, they were important when measuring sensitive
processes such as radioactive decay.
As of May 2019, the kilogram was redefined by setting the Planck constant to be
6.62607015 × 10−34 kg m2 s-1 .
Regardless of how the kilogram is defined, you might be wondering why we can use a
beam balance to compare masses. We can use a beam balance because weight is proportional
to mass according to the weak equivalence principle, which says that gravitational mass
is equal to inertial mass. Gravitational mass is the mass that determines gravitational
forces between objects; whereas the inertial mass determines the acceleration of an object
experiencing a given force. At the time of this writing, experiments show that the two
masses are equivalent to about one part in 1012 , and future experiments are planned for
even more accurate testing.
1.2.2.2 Force
Like mass, the concept of force is one with which we are intuitively familiar. A force is
essentially a push or a pull in a particular direction. You push forward on a shopping cart
to move it in the direction you wish to go. However, if someone steps in front of your cart as
it is moving, you would pull back on the cart to make it stop. From our everyday experience,
we know that forces cause motion and that multiple forces can be acting on an object. For
example, if you are in the gym lifting weights, then you are exerting a force to raise a
dumbbell. If you let go of the dumbbell, it will fall, demonstrating that the force of gravity
is also acting on the dumbbell. Hence, we see that in order to understand the motion of an
object, we need to know about all of the forces acting on the object. When accounting for
all of the forces, we need to know both the magnitude (or amount) and direction of each
force. Similar to displacement, velocity, and acceleration, force is a vector quantity.
i i
i i
i i
The unit of force is the newton (abbreviated N) and 1 N is the total amount of force
needed to provide an acceleration of 1 m/s2 to a 1 kg mass. We also know from everyday
experiences that there is a direct linear relationship between force and acceleration. For
example, a 4 N net force will cause a 1 kg object to accelerate 4 m/s2 . Of course, the
acceleration is caused by the vector sum of the forces. When you hold your cell phone, you
are exerting a force upwards that matches the downward force of gravity, hence the two
forces are equal in magnitude but opposite in direction, and therefore, their sum is zero.
Hence, the cell phone does not accelerate.
We have now laid out all of the tools needed to describe and explain an object’s motion.
Next, we will discuss the core topic of classical mechanics, Newton’s Laws of Motion, which
will explain the role of force and mass in determining the motion of an object.
A particle’s velocity remains constant if the net force acting on the particle is
equal to zero.
Let us look carefully at what the first law says. The first law says that a particle’s
velocity remains constant if the vector sum of the forces (net force) is zero. The term “net
force” is important because there can be forces acting on the particle, but if all of those
forces sum to zero, then the particle’s velocity doesn’t change. The first law is sometimes
referred to as the law of inertia because it says that an object will continue moving with a
constant velocity forever, if there are no net forces acting on it. Inertia is, simply put, an
object’s resistance to acceleration. Hence, the particle’s inertia will ensure that the particle
will continue moving in a straight line with a constant velocity or remain at rest (zero
velocity), until it experiences a force that will change the particle’s speed, direction, or
both.
i i
i i
i i
We know that acceleration is the change of velocity and hence, the first law states that
a nonzero net force is needed to cause an acceleration.
Finally, the first law provides a definition for the term equilibrium. A particle is in
equilibrium if the acceleration of the particle is equal to zero. Hence, one condition for
equilibrium is that the net force acting on a particle must be zero. We will later see that
when we study rotational dynamics, we will also need the net torque to be equal to zero as
an additional condition for equilibrium.
Before moving on to Newton’s second law, it is useful to compare our statement of the
first law to the more colloquial “An object in motion remains in motion, and an object
at rest stays at rest unless acted upon by an unbalanced force.” While there is nothing
technically wrong with this statement, notice that it is longer than our statement of the law
on the previous page. The colloquial version breaks down motion and rest as two different
behaviors. However, our use of the word velocity takes both states (motion and rest) into
account because being at rest simply means that the particle’s velocity is equal to zero. As
mentioned previously, the term motion is not one of our descriptors and therefore it is not
clear what is remaining constant. Additional terms can be added to explain motion, but at
the cost of conciseness. Finally, the term unbalanced force is also not as clear as the term
net force. The term net force means vector sum of forces, a clearly defined term. While the
colloquial statement of the First Law may use words familiar to most people, a physicist
would prefer a statement similar to the one in the box above, due to its use of precise
language.
A particle’s time rate of change of linear momentum, p, is equal to the net force, F,
applied to the particle:
F = ṗ (1.3.1)
where the linear momentum of a particle with a mass m and velocity v is defined to be:
Linear Momentum of a Single Particle
p = mv (1.3.2)
Equation (1.3.1) might not be the way you are used to seeing Newton’s second law. If
we consider a system with constant mass m such as a single particle experiencing a nonzero
net force, then we have:
dp d
F= = (mv) (1.3.3)
dt dt
F = ma (1.3.4)
and we recover the more familiar form of the second law. However, when studying the
motion of an object whose mass m is changing with time, such as in the case of a rocket,
i i
i i
i i
the form of the second law presented in (1.3.1) will be the one needed in order to derive the
object’s equations of motion.
Although we will use the form of Newton’s second law F = ma most often for calcula-
tional purposes in this book, we find it helpful to think of the second law in the form:
F
a= (1.3.5)
m
We find this form to be more useful when understanding the concepts behind Newton’s
second law. Equation (1.3.5) tells us:
• Acceleration is in the same direction as the net force. There is no minus sign in front
of F to denote an acceleration in the opposite direction, nor is there any mathematical
transformation done to the vector F to rotate it.
• The magnitude of the acceleration is directly proportional to the magnitude of the net
force. In other words, net forces with a large magnitude produce larger magnitude
accelerations than net forces with small magnitudes. Mathematically, this can be seen
because F appears in the numerator of the fraction in the right-hand side of (1.3.5)
and by thinking of the equation as: |a| = |F|/m.
• Mass “resists” acceleration. The mass m appears in the denominator of the fraction
in the right-hand side of (1.3.5). Large denominators result in smaller overall fractions
when compared to a small denominator with the same numerator (i.e., 1/4 < 1/2). In
other words, with the same given force, objects with a larger mass experience a smaller
acceleration, and smaller mass objects experience a larger acceleration.
All three of the above bullet points are contained in the one simple equation (1.3.5)! This is
one of the reasons why physicists prefer math as the language for describing the universe.
A lot can be said in one simple equation. As a physicist, you should learn how to “read
equations” like we did above.
Notice the consistency between the first and second laws. If F = 0, i.e., a zero net force
is acting on the particle, then a = 0, and the particle’s velocity is not changing, just as
stated in the first law. Likewise, we could have restated Newton’s first law as: “A particle’s
momentum remains constant if no external net force acts on the particle.”
As we mentioned, the second law will be used to produce differential equations which
describe the motion of a particle. As a simple example, we will consider a particle moving
in one dimension (along a line) under the influence of a constant force, where a = ẍî and
F = F î. Recall that the dots above the x denote a second derivative with respect to time.
We can use (1.3.5) to write:
ẍ = F/m = a (1.3.6)
Equation (1.3.6) says that the solution x(t) is a function such that its second derivative is
equal to the constant, a. Recall from simple algebra that an equation like x + 3 = 10 says x is
the number that when added to 3 gives 10. In that case, the solution, x, is a number (x = 7).
The second-order differential equation (1.3.6) is similar to our algebra problem (but much
harder!), except instead of finding a number, you are asked to find a function. Differential
equations are the mathematical language used by physicists to describe the motion of a
particle. Be prepared to solve them both in closed form (analytically) and numerically. Just
to whet your appetite, if a is constant, the solution to (1.3.6) is:
i i
i i
i i
1
x(t) = x0 + v0 t + at2 (1.3.7)
2
where x0 = x(0) and v0 = v(0), are the initial position and velocity, respectively. You can
check this for yourself by taking two derivatives of (1.3.7) and showing that it satisfies
(1.3.6) in the case of a constant a. In Chapter 2, we will go through the steps to derive
(1.3.7) and other solutions to differential equations. If you aren’t certain how we obtained
(1.3.7) from (1.3.6), don’t worry, you should understand it by the end of Chapter 2.
Finally, we can return to the weak equivalence principle. Suppose a particle is near the
surface of the Earth and experiences only the force of gravity (weight W ):
F =W
minertial a =mgravitational g (1.3.8)
where we have dropped the vector notation, assumed “downward” is the positive direction,
and used g for the acceleration due to gravity (9.8 m/s2 ). The mass minertial in this equation
is the m from Newton’s second law, i.e., it is the mass that “resists” acceleration, and we
call this the inertial mass. The mass mgravitational is the mass that is affected by gravity
and is called the gravitational mass. The weak principle of equivalence says that minertial =
mgravitational , and therefore the masses cancel, and a = g for a freely falling body near the
surface of the Earth. If the weak equivalence principle were not true, then the mass of an
object would affect its acceleration due to gravity. So far, physicists have not been able to
detect any mass dependence (outside experimental error) on an object’s acceleration due to
gravity, even in the most sensitive of experiments.
If object 2 exerts a force F12 on object 1, then object 1 exerts a force F21 on object 2
such that:
Notice that the minus sign and lack of a scalar multiple in (1.3.9) says that the force
exerted by object 2 on object 1 is equal in magnitude (lack of scalar multiple which, if
present, would change the magnitude), but opposite in direction (denoted by the the minus
sign) to the force object 1 applies on 2. Many of the interaction forces we will study in this
book will be central forces. Central forces are forces that act along the line that joins the
centers of the two interacting objects, such as gravity and electrostatic forces, and obey
Newton’s third law. Velocity-dependent forces are not central, and the third law may not
apply. An example is the force between two moving electric charges; the Lorentz force is
i i
i i
i i
velocity-dependent, and the magnetic force vectors between the two charged particles do
not lie along the same line, hence the resulting net forces do not obey Newton’s third law.
Let us take (1.3.9) a little further:
Not surprisingly, we see that the accelerations of each object are in opposite directions. If
we continue manipulating (1.3.11), we find:
m1 a2
=− (1.3.12)
m2 a1
where we have taken the magnitude of the acceleration vectors. Notice that ratio of the
accelerations is inverse to the ratio of the masses. In other words, if m1 > m2 then a1 < a2 ,
in order for the right-hand side of (1.3.12) to be a fraction greater than one. As an example,
consider the classic scenario of a mosquito hitting the windshield of a moving automobile.
Let the mosquito be object 2 and the car be object 1. Clearly, the automobile has more mass
than the mosquito, m1 > m2 , and according to (1.3.12), the acceleration of the mosquito is
greater than that of the automobile, a2 > a1 . It is a bad day for the mosquito. Another way
to interpret (1.3.12) is that an object cannot accelerate without another object accelerating
in the opposite direction.
Next, we need to address one other important fundamental concept for classical mechan-
ics: reference frames. Once we have a working understanding of reference frames, we will
have finished laying out the foundations of classical mechanics, and we will be ready to do
some physics!
i i
i i
i i
y y'
xS'
xS
x x'
S S'
Figure 1.1: The frame S 0 is moving at a constant velocity with respect to the fixed frame,
S. The particle (black dot) is measured to have a position xS in the reference frame S and
xS 0 in the reference frame S 0 . The distance between the references frames is d.
xS = xS 0 + d (1.4.1)
allows the two frames to consistently describe the motion in each frame. First, we will
consider the case where S 0 is moving at a constant speed relative to S, in other words, u = d˙
is constant. Differentiating (1.4.1) with respect to time will give us the transformations of
the velocity between coordinate systems:
Hence, if the observers in each frame wanted to determine if they are consistently measuring
the velocity of the particle, then they can insert their measured velocities into (1.4.3) in
order to check the other’s result. Next, we compute the acceleration transformation:
where we used the fact that u (the relative velocity between frames) is constant. Note that
(1.4.5) says that the measured accelerations in the two frames are the same, thus (1.4.5)
i i
i i
i i
also tells us that the forces measured on the particle in each reference frame are the same
and therefore Newton’s laws hold in their typical form in each frame. In other words, if we
measured the acceleration of the particle in each reference frame, we would find that the
measured acceleration can be accounted for by considering all of the forces acting on the
particle. The reference frame, S 0 is called an inertial reference frame because it is moving
with a constant velocity. Newton’s second law holds in inertial reference frames.
Let us now contrast that with the case of S 0 accelerating with respect to S. In this case,
(1.4.3) doesn’t change, and (1.4.5) becomes:
aS = aS 0 + u̇ (1.4.6)
and u̇ 6= 0. Now the two measured accelerations are not the same, and therefore the measured
forces in each frame are different. In particular, the acceleration of a particle in a noninertial
reference frame cannot be accounted for by summing all of the forces acting on the particle.
For example, on certain carnival rides that involve rotation, you will experience an apparent
outward force which is not caused by any forces acting on you from the ride itself. This is
an example of a noninertial reference frame where Newton’s laws no longer hold in their
standard form because there is an acceleration measured in S that is not apparent in S 0 . As
we will see in a later chapter, we can modify Newton’s laws in order to address the case of
noninertial frames. Such modifications involve treating the acceleration u̇ as coming from
an inertial force, a force that is not created by physical interactions, but rather is due to an
accelerating frame. While the surface of the Earth typically approximates an inertial frame,
the Earth does rotate and revolve around the Sun; hence, the velocity of a reference frame
“glued” to the Earth’s surface is not constant. There are cases where the noninertial nature
of the “glued” frame needs to be accounted for. Such cases include long-distance motion
such as missile trajectories and the motion of wind and water currents.
i i
i i
i i
where we matched the lower and upper limits on each side. Notice that the lower-limit on
the right-hand side is t = 0, and the lower limit on the left-hand side is the value of v at
t = 0, similar for the upper limits. In addition, we included a prime on our variables of
integration to distinguish them from the limits of integration. Finally we integrate, noting
that a is constant:
v(t) = v0 + at (1.5.3)
The result is the velocity as a function of time, as requested by the problem.
Example 1.1 is an example of a problem that is “done by hand,” meaning that we were
able to perform the necessary mathematical manipulations to solve the differential equation
(1.5.1) without the aid of a computer. Equation (1.5.3) is called a “closed form” solution to
the differential equation (1.5.1) because it gives a specific solution consisting of functions
and mathematical operations. What is considered as “closed-form” is somewhat arbitrary
because, for example, a solution in the form of an infinite sum may not be considered in
“closed-form.”
As you will see in this book, solving physics problems very often involves finding the
solution to differential equations. Some of those equations have no closed-form solution,
while others are very difficult to find. In those cases, computation can be extremely helpful.
For the purposes of this book, there are two forms of computation:
• Symbolic Computation: involves using a computer to help find closed-form solutions to
differential equations, integrals, eigenvalues, and much more. You enter the equation
i i
i i
i i
or integral you want solved, and the computer program returns a closed-form solution.
You may have seen tables of derivatives and integrals; symbolic computation is a much
more sophisticated version of those.
• Numerical Computation: (also sometimes referred to as “numerics”) provides a list of
numerical values representing a solution. For example, a numerical solution to (1.5.1)
can be thought of as a table of values with two columns, t and v(t), where each row
contains a specific time, t, and the value of v at time, t. Numerical solutions are often
best represented as a graph of ordered pairs (and sometimes triplets, depending on
the dimension of the problem being solved).
As an example of each type of computation, we will use the programming languages, Math-
ematica and Python, to solve Example 1.1 both symbolically and numerically. Again,
the point of the next two examples is to provide an illustration of each type of com-
putation. We follow each example with an explanation of what is going on in the
code. Please note that all of the code that appears in this book can be downloaded at
www.routledge.com/9781138495289
v[t]/.solution[[1]]//TraditionalForm
OUTPUT: at + v0
The code for using the interpreted-language Python to solve 1.5.1 is shown below.
v = Function ( ’ v ’ )
t = Symbol ( ’ t ’ , r e a l = True , p o s i t i v e = True )
a = Symbol ( ’ a ’ , r e a l = True )
print ( g e n e r a l s o l n . rhs )
OUTPUT
C1 + a ∗ t
The first line in the Mathematica code for Example 1.2 uses Mathematica’s DSolve com-
mand to produce a solution to the differential equation. Notice that the differential equation
and its initial condition appear as the first argument of the command (in between the curly
brackets). Notice also the “==” as opposed to “=” in the equation. Often in programming
languages, the single equal sign is used for variable assignment, whereas the double-equal
sign is used for equivalence. The first line of code stores the solution generated by DSolve in
i i
i i
i i
the variable solution (notice the single equal sign). The second argument v tells Mathemat-
ica that v is the function for which we are solving, and the third argument identifies t as the
independent variable. The second line in the code outputs the solution. The command Tra-
ditionalForm displays the output in a more readable format; without using TraditionalForm,
Mathematica would display the output differently. Notice that Mathematica’s solution is
the same as the one we derived in Example 1.1.
The Python code for Example 1.2 is much more involved than the Mathematica code.
The first line, “from sympy import *” causes Python to import the SymPy library. Python,
by itself cannot perform symbolic manipulations. We needed to import the SymPy library
in order to expand Python’s capabilities. Libraries are written by experienced programmers,
and they include functions which can be used by any code that imports the library. In this
case, the SymPy library includes functions like dsolve which solve differential equations
in closed form. Note that this is not the same as the Mathematica command, DSolve.
Commands from one language typically cannot be used in another. In order to make a
distinction between the two languages, we will preserve the capitalization used in each
language. For example, if we write DSolve, then we are discussing Mathematica’s command
for solving differential equation because in Mathematica the D and S are capitalized in
the command. The letters D and S are not capitalized in Python and, therefore, when we
discuss solving differential equations in Python, we will write the command dsolve, as it
appears in the language of Python.
A word of warning: libraries, while useful, may not be free of bugs. It is best practice to
test any library functions you are using in your code to make sure they behave as expected.
Python libraries like SciPy, NumPy, and SymPy are well-tested and well-documented. Doc-
umentation for these and other Python libraries can be found online.
Mathematica automatically recognizes a, vo, and t as symbols to be manipulated. In
addition, the Mathematica syntax, v[t], allows Mathematica to automatically recognize v
as a function. However, Python by itself would think of a, vo, t as variables for assignment
and would not recognize v(t) as a function. Therefore, all of the symbol identifications
need to be specified in Python. The Symbol and Function commands from the SymPy
library do exactly that. The general solution of the differential equation is given by the
SymPy command dsolve, whose arguments are the differential equation and the function
for which the differential equation is being solved. Notice that in Python, the differential
equation is written such that the right-hand side is equal to zero, and only the left-hand
side of the differential equation is entered. At the time of this writing, Python’s dsolve
command cannot solve the differential equation with the initial conditions at the same time
(except for power series solutions), unlike Mathematica’s DSolve. With more code, one could
perform the substitutions, however, they are not included here in order to prevent (further)
confusion. The output of the Python program is included below the word, OUTPUT. The
word, OUTPUT, and the line below it are not part of the Python code; we included it only
to show the output of the program. In this book, we will often show the output of a Python
program in this manner. Note the the C1 in the output represents a constant of integration.
In Example 1.2, we see that Mathematica requires a much shorter code to solve this
problem than Python because Mathematica is designed, in part, for symbolic manipulations.
Computer algebra systems (CAS) like Mathematica and Maple are very convenient and
easy to use when performing symbolic computations. The disadvantage is that they are
both proprietary and more expensive than other options, such as the open-source free CAS
SageMath (https://www.sagemath.org/).
Next, we will look at an example of a numeric computation and compare the two lan-
guages in that context.
i i
i i
i i
30
25
velocity (m/s)
20
15
10
0
0.0 0.5 1.0 1.5 2.0 2.5 3.0
time (sec)
The Python code used to solve (1.5.1) is shown below. For brevity, we did not include
the resulting graph which is the same as the one above.
import numpy a s np
from s c i p y . i n t e g r a t e import o d e i n t
import m a t p l o t l i b . p y p l o t a s p l t
vo = 1
def velderiv (v , t ) :
a = 9.8
dvdt = a
r e t u r n dvdt
t i m e s = np . l i n s p a c e ( 0 , 3 , 3 0 )
v e l o c i t y = o d e i n t ( v e l d e r i v , vo , t i m e s )
plt . p l o t ( times , v e l o c i t y )
plt . x l a b e l ( ’ time ( s e c ) ’ )
plt . y l a b e l ( ’ v e l o c i t y (m/ s ) ’ )
plt . show ( )
i i
i i
i i
In order to produce a numerical solution for Example 1.3, both languages needed to
have numerical values for all variables. In this case, we needed to specify a and v0 . In the
Mathematica program, we used the first line of the code to specify the value of a, but we
defined the initial value of v in the Mathematica command NDSolve, which produces a
numerical solution to a differential equation. Notice that the arguments for NDSolve are
similar to those of DSolve, i.e., the order of the arguments are: equation to be solved and
initial values, the function for which we are solving (v), and identification of the independent
variable, t. This time, however, we needed to specify the range of t values for which v is
being solved. Recall the idea of a numerical solution being like a table of values containing
columns t and v(t), so we need to tell the computer when to start finding the solution (t = 0)
and when to stop (t = 3). The final line of the code produces a graph of the solution. Notice,
we get the expected line with y-intercept of 1.0 and slope of 9.8. Everything after {t, 0, 3}
(which dictates which values of t should be used to make the plot) in the final line of the
code are formatting commands which only affect the visual appearance of the graph. We
show the formatting commands so that you may make similar plots.
In the Python example, we needed to import three libraries, NumPy, SciPy, and Mat-
plotlib in order to perform the task of numerical integration. From the SciPy library we
imported only one function, odeint which will be used to numerically solve the differential
equation (1.5.1). Later in this book, we will talk about how to numerically solve differential
equations without using the libraries, but for now, we wanted to use functions included in
the aforementioned libraries in order to simplify the code. The NumPy and SciPy library
contain functions and algorithms critical for scientific computing. The functions contained
in NumPy and SciPy will greatly increase the speed of code written in Python, and we
strongly encourage you to use them whenever possible. The Matplotlib library is a plotting
library that will allow you to make graphs in Python. After the libraries are imported, we
defined the initial condition variable vo and the differential equation. In Section 13.2, we
comment more on libraries, the risks of using other people’s code, and the odeint command.
To solve the differential equation numerically, we created a new function in Python
that we called velderiv, which will contain the differential equation we are solving. User-
defined functions, like velderiv, are convenient in programming languages when a particular
calculation needs to be repeated many times. In Python, function definition is done using
the command def and is demonstrated in the above code. The arguments of the function
are included in parentheses following the function name. In this case, the arguments are v
and t. The next few lines contain the actual calculation of the function. Notice we included
a local variable dvdt, which stays within the function and is equal to the first derivative of
v. Equation (1.5.1) states that the first derivative of v is equal to a. The last line of the
function begins with the command return, which tells Python what value to return when
the function is called. In this case, we return the value of the variable dvdt. Hence, the
function velderiv returns the first derivative of v.
For the purposes of solving a differential equation, we write the function of the form,
dx
= f (x, t) (1.5.4)
dt
where the derivative is on the left-hand side and everything else on the right. As described
in the paragraph above, the function velderiv contains the right-hand side of the differential
equation (the side without derivatives). We will expand this procedure to second-order
differential equations in Chapter 2. After the differential equation is defined as a function,
we need to tell Python for which values of t we will be computing v. In this case, times is
an array which includes the values of t, which will be used to compute v. Arrays are lists
of values. In this case the array times contains the list [t0 , t1 , . . . , t30 ] = [0, 0.1, . . . , 3.0], and
i i
i i
i i
hence we will be solving for v(0), v(0.1), . . . , v(3.0). Next, it is time to solve the differential
equation. The command, odeint, used to solve the differential equation comes from the
scipy.integrate subpackage of the SciPy library. The odeint command requires the differential
equation (as a user-defined function, velderiv), the initial value (v0 ), and the list of times,
t as arguments. The result is an array, we called velocity, and it contains values [v0 , v1 , . . .],
where vi = v(ti ). The last four lines of the program set up the graph of v versus t, which is
not shown.
Again, the Python code is longer than the Mathematica code. Does that make Mathe-
matica better than Python? That answer depends on what you are trying to do. If you are
interested in quick development and implementation, then Mathematica might be a very
useful tool for you. However, as we mentioned before, Mathematica is proprietary and more
expensive than Python. The proprietary nature of Mathematica means that you don’t have
access to the source code, and this can be a problem if you want to know exactly what
the commands are doing. Python is open source, so if you want to know what is “going
on under the hood,” you can find out. The open source nature of Python is valuable when
doing research, and you want to identify if unexpected results are due to bugs in code or
new science. That said, Mathematica is a language of its own, and you can write your
own programs in it. Finally, Python is free, and your budget may dictate your choice. We
will elaborate more on language selection in the next section. When working on their own
research problems, the authors of this book will often use both programs, choosing the most
appropriate tool for the particular job at hand.
It should be pointed out that sometimes, a closed-form solution to a differential equation
doesn’t exist or is difficult to find. In those cases, we rely on numerical solutions to give
us the information we need. You may think that numerical solutions are of limited value,
but as we will see throughout this book, there is a lot that can be found using numerics.
In Chapter 13, we will study nonlinear oscillators, and we will learn how to obtain a lot of
information from numerical solutions.
It is easy to walk away from an undergraduate physics education thinking that all physics
problems have closed-form solutions. This conclusion arises from the types of problems that
undergraduate students solve as part of their education. The truth of the matter is that
most problems physicists work on outside of the classroom require numerical solutions. The
differential equations governing real-world systems are often complex and not solvable in
closed-form. In those cases, numerics may be the only option to gaining any kind of insight
into the problem. A physicist with strong computational skills will be well-prepared to
tackle a wide variety of problems, not just in classical mechanics but in any field of physics
or engineering.
Furthermore, physicists find themselves working on a variety of problems outside the
traditional subfields of physics. Physicists often end up working on problems in finance,
economics, biology, climate science, and materials science, just to name a few. Today’s
physicists are involved in modeling economies, disease propagation, and social networks.
These types of problems allow physicists to apply their skills, including strong computational
and modeling skills, to interesting problems. It is our hope that this book will not only
prepare you well in physics but also set you on the path of developing strong computational
skills, so that you may tackle the exciting problems both inside and outside of physics,
wherever your career may take you.
i i
i i
i i
i i
i i
i i
problem. Python, an interpreted language, has a library called SymPy, which, as we saw
above, allows Python to become a CAS.
The type of language that you use depends on the nature of the problem you are trying
to solve. If you need help simplifying the mathematical form of an equation, a CAS is prob-
ably your best choice. While SymPy may be capable of mathematical manipulations, you
may find Mathematica easier to use for that problem, since Mathematica was specifically
designed for, among other things, symbolic manipulations. Interpreted languages are help-
ful for problems that are not computationally intensive, such as numerically solving certain
differential equations. The speed of compiled languages makes them ideal for computation-
ally intensive programs, such as climate models or modeling the motion of many stars in a
galaxy.
In this book, we are not going to be tied to one specific language. The focus of this book
is on classical mechanics and computation, not on a particular programming language.
The problems we will solve in this book are not computationally intensive, i.e., they won’t
demand a lot of computer resources or time. Hence, we will not be using a compiled language
in this book. For the sake of consistency, we chose to demonstrate the computations using
two languages, Mathematica and Python. The exercises and examples done in this book
can be done in any language you prefer, and, when appropriate, the code will be written in
such a way that it will be easy to reproduce in your language of choice. However, we felt
it would be best to avoid using too many different languages in the examples done in the
text.
Our choice of Mathematica and Python as example languages is not random. Both lan-
guages are used extensively in physics, and both languages are mature. The code that we will
write will focus on core commands that have been around for a long time and are not likely to
become outdated anytime soon. Python is free to use and available online for any platform.
You’ll likely want to download a distribution like Anaconda (https://www.anaconda.com)
which comes with an integrated developer environment (IDE). Mathematica is not free,
however, there are open source alternatives, such as SageMath which is available online
(https://www.sagemath.org/). SageMath was developed as an open-source alternative to
Mathematica. While there will be syntax differences between Mathematica and SageMath,
the symbolic calculations done in this text in Mathematica can, for the most part, be done
in SageMath.
You do not need to have experience with any programming languages to understand the
material presented in this book. Such previous experience will undoubtedly help in under-
standing the code contained in the book. However, we will introduce the necessary program-
ming concepts as they are needed, and we assume no programming experience. Mathematica
has extensive and very useful help files that we recommend going through. If you’d like to get
a head start on Python, we recommend Code Academy (https://www.codecademy.com/)
and the most recent edition of [Kinder and Nelson(2015)]. A quick way of getting your
feet wet and learning some basic programming skills is through using your favorite search
engine. Searches such as: “How to solve a differential equation in Mathematica” or “How
to integrate an equation in Python” often lead to many useful pages, which will at least
give you the syntax for how to solve a specific problem. While this is no substitution for an
introductory computer science course or a course in numerical analysis, it will certainly get
you started.
i i
i i
i i
i i
i i
i i
acting agents whose collective behavior cannot be explained by the dynamics of the indi-
vidual agents alone. Examples of complex systems include power grids, global climate, and
economies. Again, these examples aren’t strictly “physics problems” but rather problems to
which physicists have contributed greatly. Progress in the fields of nonlinear and complex
systems has required people trained in physics, notably classical mechanics and statistical
mechanics.
Another important field involving classical mechanics is fluid mechanics and turbulence.
There is a famous story of Werner Heisenberg, that on his deathbed he said, “When I
meet God, I am going to ask him two questions: Why relativity? And why turbulence? I
really believe he will have an answer for the first.” Variations of this story are told about
other physicists as well. The point is that turbulence continues to be a difficult problem in
classical mechanics. Improvements in our understanding of turbulence will greatly benefit
many technologies that involve fluid flow, such as pipelines and flight.
The paragraphs above contain just a few examples of the relevance of classical mechanics
to the modern world. We have left out many other examples such as mechanical engineering,
oceanography, and space flight. The point is that ignoring classical mechanics or treating
it simply as a step to get to something else, is a mistake. A fundamental understanding
of classical systems is critical to becoming a complete physicist. In many ways, classical
mechanics serves as the foundation of the field of physics itself. A physicist needs a solid
foundation of classical mechanics in order to understand other sub fields of physics.
One of the many benefits of learning classical mechanics is that it gives you the ability
to identify what factors influence a system (physical or otherwise) and how to express those
factors mathematically. Much of this ability centers on the idea of learning how to specify
a system, identify the forces acting on the system, and then writing down the relevant
equations of motion that govern the system. It is a powerful and useful way of thinking,
that will be a critical component of your training as a physicist and is a valuable skill
regardless of how you end up using your physics training.
i i
i i
i i
Newton’s First Law: A particle’s velocity remains constant if the net force acting on
the object is zero.
Newton’s second law: A particle’s time rate of change of momentum dp/dt, is equal to
the net force F applied to the particle :
F = dp/dt = ṗ (1.7.4)
Newton’s third law: If object 2 exerts a force F12 on object 1, then object 1 exerts a force
F21 on object 2 such that:
F21 = −F12 (1.7.5)
All measurements are made with respect to a reference frame. A reference frame is a
choice of origin (spatial x and temporal t), and a set of axes with respect to which all
measurements are made.
Two reference frames S, S 0 are called inertial reference frames when Newton’s laws
hold, and the forces measured in S 0 are the same forces as those measured in S.
Two reference frames S, S 0 are called noninertial reference frames when Newton’s laws
no longer hold in their standard form, because there is an acceleration measured in S that
is not apparent in S 0 .
Computation is an important tool for physicists. In this book, we’ll focus on two
types of computation, symbolic and numerical. Symbolic computation involves using a
computer to help find closed-form solutions to differential equations, integrals, eigenvalues,
etc. Numerical computation or “numerics,” provides a list of numerical values representing
a solution to differential equations, integrals, eigenvalues, etc.
Computer programming languages can be loosely categorized as a computer alge-
bra systems, interpreted languages, and compiled languages. A computer algebra system
(CAS) is a software that can manipulate mathematical expressions, and can solve a vari-
ety of mathematical problems in closed-form. An interpreted language is one that can be
executed directly from the source code, using software called an interpreter. A compiled
language is one whose source code needs to be compiled by software called a compiler,
which transforms the source code into machine code.
i i
i i
i i
speed of the moving frame S 0 relative to the rest frame S. For example, time passes at
different rates depending on the speed of an observer, by the rule: ∆t = γ∆t0 , where
∆t is the interval of time measured in the rest frame S, and ∆t0 is the duration of
the same interval of time as measured in a moving frame S 0 . If a rocket leaves the
Earth and travels for 100 years, as measured by people on Earth, at a speed 0.75c,
how much time has passed for the travelers on the rocket?
12. In the theory of special relativity, we need to alter the transformation equations
between two frames. Let S be a rest frame, and S 0 a frame moving at a speed v
relative to S. Suppose that S 0 moves in the x direction, in other words, S 0 moves
parallel to the x-axis of the frame S. The theory of special relativity then states that
measurements made in S and S 0 are related by the formulas:
x0 = γ(x − vt)
y0 = y
z0 = z
vx
t0 = γ t − 2
c
i i
i i
i i
where γ = 1 − v 2 /c2 . Now consider that two events in S occur at two locations x = 0
p
and x = a at time t = 0. Find the times of the two events as measured in S 0 . Notice
that events that are simultaneous in S are not simultaneous in S 0 . Which event was
seen first in S 0 ?
Section 1.5 Computation in Physics
The problems in this section are different from the rest of the book. They are intended to
prepare you to better understand the computational aspects of this book. The problems
may be done in any language specified by your instructor. You may want to consider doing
some of the problems in Python and/or Mathematica so that you have some familiarity with
the code that appears throughout this text. As long as it is approved by your instructor,
online searches can be of significant help in solving these problems. Finally, these problems
are not exhaustive for the different types of problems you’ll encounter in this text. We will
introduce new algorithms and commands as they are needed.
13. Assign the values 5 and 6 to the variables, a and b. Then compute:
(a) a + b
(b) b − a
(c) ab
(d) b/a
(e) ab
14. Compute the following:
(a) sin(0.25π)
(b) cos(0.25π)
(c) e3
√
(d) 7
15. Define an array (sometimes called a list) called t which contains the numbers 0 through
2π in steps of:
(a) 0.1
(b) 0.01
(c) 0.25
16. Define the function f (x) = e−mx cos(2πkx) and compute f (x) for three different com-
binations of values of x, m, and k.
17. Plot the function from Problem 16 with k = 1 and m = 0.5 for values x starting at
x = 0 and ending at x = 10. Note that you may need to create an array for x, and
that in that case use steps of 0.01. Label the axes and choose a color other than the
default for your curve. Save the resulting graph to a .jpg or .png file.
18. Plot f (x) = x2 and g(x) = x on the same graph for the range x = 0 to x = 10. Each
function should have its own color in the graph.
i i
i i
i i
19. A conditional statement is an important tool for helping computers deal with contin-
gencies. They take the form “if-then-else.” Conditional statements allow for computers
to do different things, depending on whether or not a condition is true. In other words,
if a condition is true, then do one thing, if it is not true (else) then do something dif-
ferent. For example, if a variable x is less than 5 assign the value 1 to the variable a;
otherwise, assign the value 0 to a. Write a code in which you assign a value to the
variable x and that prints “x is less than 5” if x < 5 , and prints “x is greater than or
equal to 5” if x ≥ 5.
20. Using the conditional if statement, described in Problem 19, define the piecewise
function
(
0 x<5
f (x) =
x2 x ≥ 5
and evaluate f (x) for x = 3 and x = 7 and plot f (x) for the range x ∈ [0, 10].
21. Computers are very good at repeating tasks over and over again, in a procedure called
a loop. The basic types of loops are for loops, do loops, and while loops. Using any
kind of loop, compute 6! (factorial) and display the result.
22. Using the methods learned from the problems above, compute the first 20 values of the
Fibonacci sequence. If you are unfamiliar with the Fibonacci sequence, do an online
search. You will need to use loops, and you may want to save your results into an
array. Look up how to append values to an array and how to recall specific values for
an array in the language of your choice.
23. Use symbolic computation to solve the following equations for x. Note that you will
get complex roots for some of the solutions.
(a) 7x + 5 = 0
(b) x2 − 5x + 2 = 0
(c) x3 + 7x − 5 = 3
24. Use symbolic computation to solve for (x, y):
2x − 5y = 7
x+y = 2
25. Use symbolic computation to solve the following differential equation for the unknown
function f (x):
df
= −x2 + 3
dx
for the initial condition f (0) = 0. Plot the solution for x = 0 to 3 with a step size of
0.1.
i i
i i
i i
If you are using Python, you will want to consider the nquad command from the SciPy
integrate library.
31. Numerically evaluate the integral:
Z 2 Z 1−y
xydxdy.
0 0
As with Problem 30, you will want to consider using the nquad command if you are
using Python.
i i
i i
i i
CHAPTER 2
Single-Particle Motion in
One Dimension
In this chapter, we will examine one-dimensional motion, i.e., motion along a line. It is
sometimes the case that a particle’s motion need only to be described along one direc-
tion. Furthermore, a careful study of one-dimensional motion will be a useful foundation
for understanding more general motion in higher dimensions. In this chapter, we will give
several examples of solving Newton’s second law, F = ma in one dimension. We will con-
sider several types of forces: both constant and those which depend on time F (t), velocity
F (v), and position F (x). In addition, we will discuss and demonstrate two different uses
of computers to solve physics problems: how to use computer algebra systems (CAS) to
obtain the analytical solutions of Newton’s second law, and how to obtain numerical solu-
tions of ordinary differential equations (ODE) using software packages and by using the
Euler method.
F = ma (2.1.1)
Recall that acceleration is the first derivative of velocity v with respect to time and the
second derivative of displacement x with respect to time. To solve for the equations of
29
i i
i i
i i
motion, we will think of (2.1.1) as a differential equation, by rewriting (2.1.1) in the following
ways:
d2 x
F =m (2.1.2)
dt2
dv
F =m (2.1.3)
dt
dv
F = mv (2.1.4)
dx
Notice that each of the above equations is a differential equation which can be solved
once the net force F acting on the particle is specified. Each of the above equations yields
a different equation of motion: (2.1.2) can be solved for x(t), (2.1.3) for v(t), and (2.1.4) for
v(x). Equation (2.1.4) comes from the chain rule:
dv dv dx dv
m =m = mv (2.1.5)
dt dx dt dx
Throughout the rest of this chapter, we will solve (2.1.1) for several different cases where
forces are constant (F = F0 ), time-dependent (F = F (t)), velocity-dependent (F = F (v)),
and position-dependent (F = F (x)). However, before solving (2.1.1), we will make a few
comments about differential equations in general.
2. The ordinary derivative implies that x is a function of only one variable t, which is
the variable of differentiation.
3. The number of arbitrary constants in the general solution of (2.2.1) is equal to the
order of the ODE.
To solve (2.2.1), we needed to separate the variables of the equation. Colloquially speaking,
this means getting all the terms with x on one side of the equation and all the terms that are
either constant or depend on t on the other. This process is called separation of variables
i i
i i
i i
and is performed by treating the derivative as a fraction and multiplying both sides of
(2.2.1) by dt,
dx =7dt (2.2.2)
Z Z
dx = 7dt (2.2.3)
x =7t + c (2.2.4)
To solve (2.2.1), we carried out the integral after separating out variables. Note that
both integrals would produce constants of integration, but since both are constants, we can
combine them into one arbitrary constant. You can double-check the solution by computing
the derivative of 7t + c, to check that it satisfies (2.2.1).
What happens in the case of second-order ODEs? Second-order ODEs are common
in physics. There are many techniques to solve them, but we will demonstrate only one.
Consider the ODE,
d2 x
=7 (2.2.5)
dt2
Separation of variables does not make sense here, because we typically do not integrate
terms like d2 x. However, we can define a new variable v, such that, v = dx/dt. Then (2.2.5)
becomes,
dv
=7 (2.2.6)
dt
which we know gives the answer v(t) = 7t + c1 , where c1 is the constant of integration.
However, we want x(t), so we use v = dx/dt:
dx
=7t + c1 (2.2.7)
Z dt Z
dx = (7t + c1 ) dt (2.2.8)
where c2 is the constant of integration obtained by performing the above integral. Hence, we
pick up an additional constant of integration in our solution, giving two arbitrary constants
for the solution of the second-order ODE (2.2.5). Loosely speaking, we see that the number
of arbitrary constants in the solution of an nth -order ODE is equal to n, because we need to
do n integrations in order to solve the equation, and each integration produces an arbitrary
constant.
Next, we return to (2.2.1). Suppose (2.2.1) was an equation we wanted to use in order
to find the position of a particle as a function of time. The infinite number of solutions is
not helpful. Which solution describes the actual path taken by the particle? In order to
specify the particular solution for an ODE, we need to include initial conditions, the value
of our function at a particular time (normally at t = 0). Suppose we know that at t = 0, the
particle is at a position, x(0) = 3. Then we can solve for the arbitrary constant by inserting
the initial condition into our general solution,
i i
i i
i i
which gives c = 3. Our particular solution is then, x(t) = 7t + 3. A different initial condition
will give a different particular solution. Now suppose we wanted to find a particular solution
to (2.2.5); in that case one initial condition will not be enough because it will leave one
arbitrary constant. Hence, we will need to specify both x(t) and dx/dt at a particular time
(usually t = 0). Suppose that x(0) = 3 and v(0) = 1, where v = dx/dt. Then we have:
where (2.2.12) is the derivative of the general solution evaluated at t = 0. The result here
is that c1 = 1 and c2 = 3, and the particular solution is x(t) = 3.5t2 + t + 3. In summary, in
order to solve for the particular solution of an nth -order ODE, we need n initial conditions.
In addition, we can also specify x using knowledge of the value of x at two different times,
as opposed to knowing initial values of x and its first derivative. In classical mechanics, it is
most common to know the initial conditions of the position and velocity. However, in other
fields, such as electromagnetism and thermodynamics, it is often more common to know the
value of a function, say the temperature, at two different locations. In this case, we have
what is known as a boundary value problem, and the conditions that provide the constants
of integration are known as boundary conditions.
There are many other properties of ODEs such as linear superposition, that we will
explore in this book as we need them. For now, we know enough about ODEs to get started.
Let’s get back to the physics.
dv = adt (2.3.2)
which can be integrated to yield,
Z
v(t) = a dt = at + c1 (2.3.3)
where c1 is the constant of integration. The constant, c1 , can be found using initial conditions
of v(t0 ) = v0 , where v0 is the initial velocity at the initial time t0 . When t0 = 0, (2.3.3) gives
v(0) = c1 or c1 = v0 . Therefore, the solution to the differential equation, (2.3.1) is,
i i
i i
i i
Z v(t) Z t
0
dv = adt0 (2.3.5)
v0 t0
v(t) − v0 =a(t − t0 ) (2.3.6)
where we have introduced primes to the variables of integration in order to distinguish them
from the limits of integration. Notice that the lower limit in the left-hand side of (2.3.5)
corresponds to the value of v(t) when t = t0 , the lower limit of the right-hand side of (2.3.5),
and with similar considerations for the upper limits. It is very important that the limits
match on both sides of the equation.
Next, we can get an equation for x(t) by writing v = dx/dt :
dx
= at + v0 , (2.3.7)
dt
and separating variables we obtain:
Z x(t) Z t
dx0 = (at0 + v0 )dt0 (2.3.8)
x0 t0
Notice how all of the time-dependent and constant terms are on the same side of the
equation. If we had subtracted v0 from both sides of the equation before integrating, we
would have dx/dt − v0 , which does not make sense because v0 is meaningless without a
R
dv F0
v= (2.3.10)
dx m
vdv =adx (2.3.11)
Z v Z x
0 0
v dv =a dx0 (2.3.12)
v0 x0
v − v02 =2a(x − x0 )
2
(2.3.13)
where a = F0 /m was used in (2.3.11). We could have also obtained this result by eliminating
t between equations (2.3.3) and (2.3.9). The box below summarizes all of the constant force
equations, sometimes called the kinematic equations. These are equations that you should
memorize.
v(t) = vo + at (2.3.14)
i i
i i
i i
1
x(t) = xo + vo t + at2 (2.3.15)
2
In the case of a freely falling particle near the surface of the Earth, we use a = −g =
−9.8 m/s2 (assuming “down” is in the negative direction), where g is the acceleration due
to gravity. The kinematic equations become:
We used y as the position variable, which is common when describing vertical motion.
Notice that these are not different equations from the kinematic equations but simply the
kinematic equations with a specific value of a. We now look at some well-known examples
of situations in which the acceleration of the system is constant.
i i
i i
i i
Once a has been found, we can find the position x(t) and velocity v(t) of the masses
using (2.3.15) and (2.3.14). Note that for the position and velocity of m1 , we would need
to insert −a in the kinematic equations.
a a
T
m1
T
w m2
1
w
2
Figure 2.1: The free body diagrams of the two hanging masses of the Atwood machine.
To find the tension in the string, substitute (2.3.21) into −m1 a = m1 g − T , which
results in:
2m1 m2
T= g (2.3.22)
m1 + m2
Of course we could have also substituted (2.3.21) into the equation describing the forces
acting on m2 , and we have obtained the same result for T .
i i
i i
i i
mgsin( )
mgcos( )
The normal force here represents the force applied by the plane against the object
(and vice versa). The magnitude of the normal force N can be calculated from the free
body diagram, as shown in Figure 2.2, to be N = mg cos θ, where g is the acceleration due
to gravity, and θ is the angle of the inclined surface measured from the horizontal. The
component of the weight W = mg, which acts along the direction of the plane, is again
found from the force diagram to be F = mg sin θ.
The total force along the x-axis (downward direction along the plane) is:
dv F (t)
= (2.4.1)
dt m
and by separating variables and integrating from t0 = t0 to t we find :
Zt
1
v(t) − v(t = 0) = F (t0 ) dt0 (2.4.2)
m
t0 =t0
1 t
Z
v(t) = v0 + F (t0 ) dt0 (2.4.3)
m t0 =0
Once we know v(t), we can find x(t) by integrating v = dx/dt and using the initial
conditions x(t0 ) = x0 to obtain:
i i
i i
i i
Z t
x(t) − x0 = v(t0 ) dt0 (2.4.4)
t0 =0
Next, we will demonstrate an example of how to find the equations of motion for a
particle experiencing a time-dependent force. We will demonstrate the solution both by
hand and using the CAS Mathematica. Refer to Chapter 1 for comments on how to use
Mathematica to solve differential equations in closed form.
1 t 1 t
Z Z
0 0 F0 F0
v(t) = v0 + F (t ) dt = dt = tan−1 (t) (2.4.5)
m t0 =0 m t0 =0 t2 + 1 m
Does our v(t) make sense? To check, we ask what happens when t → ∞? In the limit
of t → ∞, F → 0, and we would expect that v becomes constant. In our particular case,
as t → ∞, tan−1 (t) → π/2, and therefore, our solution v(t) follows our expectation of the
velocity approaching a constant. Next, we find x(t) using (2.4.4) with x0 = 0:
Z t
x(t) − x0 = v(t0 ) dt0 (2.4.6)
t0 =0
F0 t
Z
= tan−1 t0 dt0 (2.4.7)
m t0 =0
1
F0
= t tan−1 (t) − ln t2 + 1 (2.4.8)
m 2
i i
i i
i i
hn o i
ODEsolution = DSolve m ∗ x00 [t] == F0
t∧ 2+1 , x[0] == 0, x0 [0] == 0 , x[t], t ;
position = x[t]/.ODEsolution[[1]]
F0(2tArcTan[t]−Log[1+t2 ])
OUTPUT: 2m
velocity = D[position, t]
F0ArcTan[t]
OUTPUT: m
F0 = 1; m = 1;
12 1.4
10 1.2
v, m/s
8
m
6 1.0
4 0.8
2 0.6
0 0.4
0 2 4 6 8 10 0 2 4 6 8 10
time, s time, s
Notice that sometimes curly brackets (i.e., { }) are used in Mathematica, while other
times square brackets (i.e., [ ]) are used. In Mathematica, square brackets are used to denote
the argument of a function or command. For example, we would write x(t) as a handwritten
solution to a differential equation. The notation x(t) uses parentheses to denote the fact
that the variable t is the argument of the function x. In Mathematica, we would write the
function as x[t]. Notice, we used the notation x[t] in the command DSolve in order to tell
Mathematica that x is a function of t. Also notice that the arguments of commands are
also placed in square brackets. The command D has the arguments position and t which
are inside square brackets. The D command takes a derivative of the first argument (in
this case, the equation stored in the variable position) with respect to the argument in
the second position (in this case, t). Curly brackets in Mathematica are generally used
to denote a range. Notice in the Plot command the second argument is {t, 0, 10}, which is
telling Mathematica to plot the function stored in the variable position as a function of time
(the t in the first argument of the curly brackets) starting at t = 0 (the second argument
in the curly brackets) and ending at t = 10 (the third argument in the curly brackets).
Different programming languages use square brackets, curly brackets, and parentheses in
different ways. Mathematica uses parentheses mainly to group elements together (as you
would when handwriting mathematics), however Python uses parentheses for many of the
same things for which Mathematica uses square and curly brackets. Once one knows several
programming languages, it is common to get the usage of parentheses, curly brackets, and
square brackets mixed up!
i i
i i
i i
f = −f (v)v̂ (2.5.1)
where the unit vector v̂ = v/v is along the direction of the object’s velocity (or motion) and
the minus sign is included to explicitly show that the direction of the air resistance force
is opposite of the direction of the velocity. The magnitude of the air resistance f (v) is a
function of velocity which generally takes the form,
f (v) = bv + cv 2 (2.5.2)
and therefore has both a linear and a quadratic component in v. For spheres in air the
values of the constants b and c are [Taylor(2005)]:
where D is the diameter of the sphere in meters. Fortunately, one usually doesn’t need to
include both the linear and quadratic terms for air resistance. In order to determine which
term, linear or quadratic, is the most important in a given situation, one can compute the
ratio:
cv 2 0.25D2 v 2
γ= = = 1.6 × 103 Dv (2.5.4)
bv 1.6 × 10−4 Dv
Note that v should be in meters per second, and D in meters. If γ >> 1, then the quadratic
term dominates, and the linear term can be neglected. If γ << 1, then the linear term
dominates, and the quadratic term can be neglected. However, if γ ≈ 1, then both the
linear and quadratic terms need to be included. So for a 0.22 m soccer ball, the quadratic
term dominates for speeds approximately greater than 0.003 m/s (or 3 mm/s), while for
lower speeds the linear term dominates. Both linear and quadratic terms would need to be
included in the air resistance formula if the soccer ball is traveling near 3 mm/s.
To obtain equations for position and velocity as functions of time, we will consider a
generic form for F (v). Newton’s second law (2.1.1) gives:
i i
i i
i i
dv F (v)
= (2.5.5)
dt m
Rearranging the above equation gives:
dv
dt = m (2.5.6)
F (v)
and by integrating from t0 = t0 to t , and using the initial conditions v(t0 ) = v0 we find :
v
dv 0
Z
t − t0 = m (2.5.7)
v 0 =v0 F (v 0 )
By solving (2.5.7), we can find v(t), the velocity as a function of time. Once we know
v(t), we can again use (2.4.4) to find x(t) with the initial condition x(t0 ) = x0 .
We can follow a similar procedure for finding v(x). Newton’s Second Law gives:
dv
mv = F (v) (2.5.8)
dx
After rearranging, we obtain:
v
v 0 dv 0
Z
x − x0 = m (2.5.9)
v 0 =v0 F (v 0 )
which gives:
mv0 bt
x(t) = x0 + 1 − e− m (2.5.13)
b
Here is the solution using Mathematica:
i i
i i
i i
ODEsolution = DSolve[{m ∗ x00 [t] == −b ∗ x0 [t], x[0] == xo, x0 [0] == vo}, x[t], t];
velocity = D[position, t]
bt
OUTPUT: e− m vo
vo = 1; xo = 0; m = 1; b = 1;
1.0 1.0
0.8 0.8
v, m/s
0.6 0.6
x, m
0.4 0.4
0.2 0.2
0.0 0.0
0 2 4 6 8 10 0 2 4 6 8 10
time, s time, s
The Collect command gathers together terms that involve the same powers of the
objects in the curly brackets (second argument of Collect). It is a useful command for
simplifying algebraic terms and for identifying coefficients.
Note that we have factored out the constant b from the denominator of the integral
before performing the integration.
Solving for v(t) we obtain:
i i
i i
i i
mg mg − bt
v(t) = + v0 − e m (2.5.16)
b b
We also find y(t) from (2.4.4):
Z t Z t
0 0 mg mg − bt0
y(t) − y0 = v(t ) dt = + v0 − e m dt0 (2.5.17)
t0 =t0 t0 =t0 b b
Figure 2.3 shows the plots of x(t) and v(t) found in Example 2.5. Notice that both plots
show the velocity of the object coming to a constant value as the object falls for a long
period of time.
1200 50
1000 40
800
v, m/s
30
x, m
600
20
400
200 10
0 0
0 5 10 15 20 25 30 0 5 10 15 20 25 30
time, s t, s
Figure 2.3: The plots of x(t) and v(t) from Example 2.5 using Mathematica and the numer-
ical values m = 1.0 kg, b = 0.2 Ns/m, x0 = 0 m and v0 = 0 m/s.
i i
i i
i i
For simplicity’s sake, we will denote vterminal as vt . Next, we use (2.5.7) to find v(t) by
integrating from the initial velocity v0 to the final velocity vf :
vf
dv 0
Z
t =m
−cv 02 + mg
Z0
m vf dv 0 (2.5.21)
=
c 0 vt2 − v 02
m vf
= tanh−1 ( )
cvt vt
where we used integral tables to perform the integration (Note that a CAS could be used
as well, see the code below). We can then solve for vf to yield:
h cv i
t
vf (t) = vt tanh t (2.5.22)
m
We also find x(t) from (2.4.4):
Z t Z t h cv i
0 0 t 0
x(t) − x0 = v(t ) dt = vt tanh t dt0 (2.5.23)
0 0 m
which with x0 = 0 gives:
m cv
t
x(t) =
ln cosh t (2.5.24)
c m
In the code below, we used Mathematica to perform some of the calculus needed to
solve the problem.
h h ii
integral = Assuming (vf > 0)&&(vt > 0)&&(vf < vt), Integrate vt∧ 2−
1
v ∧ 2 , {v, 0, vf}
ArcTanh[ vt
vf
]
OUTPUT: vt
velocity = vt*Tanh ;
c∗t∗vt
m
position = Integrate[velocity, t]
mLog[Cosh[ ctvt
m ]]
OUTPUT: c
i i
i i
i i
7
30
6
25
5
20
v, m/s
m 4
15
3
10 2
5 1
0 0
0 1 2 3 4 5 0 1 2 3 4 5
time, s time, s
In Example 2.6, we used the command Integrate to perform the integrals. When you
perform integrals you may not always think about the signs of variables. For example
you know that the terminal velocity is positive. However, Mathematica (and other CAS
programs) makes no such assumptions. In order to perform the integral, and get the above
result, we needed to tell Mathematica that both the final velocity (vf in the code), and
terminal velocity, (vt in the code) are positive, and that vf < vt. Physically, this is making
the assumptions: m, g, k, t, v > 0 and mg > kv 2 (i.e., the air resistance cannot exceed the
weight W = mg). The symbol && represents the logical operator AND. These assumptions
are provided in the Assuming command in the above code. So by placing the command
Integrate inside the Assuming command, we are telling Mathematica to integrate with the
listed assumptions. Notice that in the first line of code, the command Integrate is used with a
second argument {v, 0, vf }, while the position calculation has the simple second argument of
t. In the first case, we are asking Mathematica to perform a definite integral with variable
of integration, v, and with a lower limit of v = 0 and an upper limit of v = vf . In the
second case, we are asking Mathematica to perform an indefinite integral with variable of
integration t.
A comparison between the velocity of the falling body in the cases of linear air resistance
F (v) = −bv + mg and quadratic air resistance F (v) = −cv 2 + mg is shown in Figure 2.4,
where m = 1.0 kg, b = 0.2 Ns/m, and c = 0.2 Ns2 /m. The velocity on the y-axis has been
normalized by dividing with the corresponding terminal velocity. Notice that the quadratic
air resistance (dashed line) leads to the object obtaining terminal velocity in a significantly
shorter time than linear air resistance (solid line). This is not surprising; the magnitude of
the quadratic drag force will be higher than the linear force for a given velocity, v.
1.0
0.8
velocity, m/s
0.6
Out[ ]=
0.4 Linear
Quadratic
0.2
0.0
0 5 10 15 20
time
Figure 2.4: Comparison of the velocity of falling body in the presence of linear and quadratic
air resistance. The two curves have been scaled, so that they produce the same result as
t → ∞.
i i
i i
i i
1 1
Z x Z v
F (x0 )dx0 = mv 0 dv 0 = mv 2 − mv02 (2.6.2)
x0 v0 2 2
Solving for the velocity v(x) as a function of the position x:
sZ
2 x
1
v (x) = F (x0 )dx0 + mv02 (2.6.3)
m x0 2
By substituting v = dx/dt, separating the variables, and integrating with the initial
condition x = x0 when t = t0 , we obtain the relationship between position x and time t:
x
dx0
Z
m
t − t0 = (2.6.4)
2
qR
x
x0 F (x )dx + 2 mv0
x0 0 1 2
After performing the integral in equation (2.6.4), one would then try to invert this
formula to find x(t), something that is not always easy or possible to do. Often, one relies
on computer methods in such cases.
The next example will involve one of the most important problems in all of classical
mechanics, simple harmonic motion (SHM). In fact, we will devote a whole chapter to it! As
we will see in a later chapter, SHM is a common model used for small amplitude oscillations.
Although SHM is often thought of as a “mass on a spring,” it is also a useful model for
other small amplitude oscillations including pendulums, small amplitude water waves, and
loudspeakers—to name a few. Besides presenting the SHM, the next example also includes
a demonstration of using Python to find and plot the solution to a differential equation
in closed form. Python’s SymPy package was used in Chapter 1, to solve a first-order
differential equation. Here we will demonstrate the solution of a second-order differential
equation using SymPy.
i i
i i
i i
Solution:
In this case (2.6.3) becomes:
s Z
2 x
1 2 1
r
v= kx dx + mv0 =
2 k x20 − x2 (2.6.5)
− 0 0
m x0 2 m 2
which gives the velocity v(x) as a function of the position x. Next we will find x(t) by
using (2.6.4):
Z x x r
dx0
r
m −1 x m −1 x π
t= = sin = sin − (2.6.6)
2
q
k x0 k x0
m (x0 − x )
x0 k 2 02
x0
i i
i i
i i
x = Function ( ’ x ’ )
t = Symbol ( ” t ” , r e a l=True , p o s i t i v e=True )
k = Symbol ( ” k ” , r e a l=True , p o s i t i v e=True )
m = Symbol ( ”m” , r e a l=True , p o s i t i v e=True )
d i f f e q = Eq (m∗x ( t ) . d i f f ( t , t ) + k∗x ( t ) , 0 )
equation = s o l u t i o n . rhs
C1 , C2 = symbols ( ’ C1 , C2 ’ )
xo = symbols ( ’ xo ’ )
i c 1 = Eq ( e q u a t i o n . s ub s ( t , 0 ) , xo )
i c 2 = Eq ( e q u a t i o n . d i f f ( t ) . s u b s ( t , 0 ) , 0 )
s o l u t i o n i c s = s o l v e ( [ i c 1 , i c 2 ] , ( C1 , C2 ) )
print ( solution ics )
OUTPUT: {C1 : 0 , C2 : xo }
p a r t i c u l a r = s i m p l i f y ( e q u a t i o n . su bs ( s o l u t i o n i c s ) )
print ( particular )
OUTPUT: xo ∗ c o s ( s q r t ( k ) ∗ t / s q r t (m) )
i i
i i
i i
We didn’t need to specify that t, m, and k are real and positive. If we had only defined
those variables as symbols, Python would have still solved the differential equation. However,
Python would produce a cosh function instead of the cos function. The argument of the cosh
function produced has a negative sign under a radical. Using the relationship, cosh(ix) =
cos(x) we would get the same answer we obtained by hand and by specifying the nature of
the variables. It is not unusual to have to do additional algebra on results provided by CAS
algorithms, in order to get a result that is useful. This is an additional reason not to forget
your math skills and not to rely solely on the computer!
Not all differential equations can be easily solved in closed form and some cannot be
solved in closed form at all. In cases where closed-form solutions are not possible to obtain,
we can solve the differential equation using a numerical method. In the next section, we will
discuss a simple method of numerically solving ordinary differential equations.
dx x (t + dt) − x (t)
≈ . (2.7.2)
dt dt
If the value of dt is small enough, the right-hand side of the above equation will give a
good approximation of the derivative. We can then think of solving the differential equation
by finding the current value of x = x(t + dt) by using past values, x(t). This process is called
Euler’s method for solving differential equations. By inserting (2.7.2) into (2.7.1), we obtain
the formula for Euler’s method:
We can interpret (2.7.3) as using the tangent line at the current value of x, x(t), in order
to approximate the next value of x, x(t + dt). Over a short distance dt, we can expect this
local linear approximation to be a good method of finding x(t + dt) from x(t).
To implement Euler’s algorithm, follow these steps in your program of choice:
1. Define dt to be a value for the time interval small enough to produce the correct
solution.
2. Define an empty array x for the function, which is the solution to the ODE. Some
programming languages require you to define the length, or number of elements in
the array, when you define the array. If that is the case, define the length of x to be
i i
i i
i i
equal to tmax /dt (rounded up to the nearest integer), where tmax is the largest value
of t for which you want to know x(t).
3. Loop over an index i in order to enter values in the array using:
The ith element of the array x, denoted by x[i], contains the value of the function x at
t = idt. How do you know if the value of dt that you chose gives a good solution? It is
sometimes helpful to run your code with a first guess of dt, and then repeat for a smaller
value of dt. If the solution doesn’t change using the new value of dt, then that is a good sign
that you have a good value of dt and an accurate representation of the solution. Of course
it is not a proof, but the method can work well as you try to get a sense of the solution
to a differential equation. It is worth mentioning here that the errors in numerical methods
such as the Euler method are well known, and although we will not go into them here, they
are covered in typical numerical analysis texts such as [Press et al. (2007)].
You might wonder, why not just use a very small value, say, dt = 0.000001? One reason
is that if you did that, then to find the solution to t = 1, the loop would take 106 iterations!
However, a more important reason is that if dt is very small, then global truncation errors
become an issue for your solution. Local truncation error is the error that occurs at each
iteration of a numerical method. Global truncation error is the cumulative error caused by
many iterations. Hence, small values of dt both increase computing time and can introduce
more error. Furthermore, the Euler method can be unstable, meaning that the numerical
solution can grow very large when the exact solution does not. If you cannot obtain a
solution with reasonable accuracy using Euler’s method, then you will need to investigate
higher-order numerical solvers such as the four-order Runge-Kutta method, which we will
discuss in a later chapter. Even more sophisticated methods involve an adaptive value of
dt.
You might be wondering how to solve second-order ordinary differential equations
(ODE), which are more common in mechanics due to Newton’s second law. To solve a
second-order ODE, we break it up into two first-order ODEs and use the Euler method for
each one. Consider the ODE resulting from Newton’s second law in the previous example
of linear air resistance:
ma = − bv + mg (2.7.5)
d2 x dx
m =−b + mg (2.7.6)
dt2 dt
where we have written a = d2 x/dt2 and v = dx/dt. To create two first-order ODEs from
(2.7.6), we treat the velocity as a separate variable:
dx
=v (2.7.7)
dt
then (2.7.6) becomes:
dv b
=g− v (2.7.8)
dt m
The resulting equations for the Euler method are therefore:
x(t + dt) =x(t) + dt v(t)
(2.7.9)
k
v(t + dt) =v(t) + dt g − v(t)
m
i i
i i
i i
In Algorithm 2, we use Python to write a Euler’s method algorithm which solves Example
2.5. In the next few paragraphs, we will explain Algorithm 2 line-by-line.
The Python library NumPy is imported and used for convenience, although it is not
necessary. We also import and use the Python library Matplotlib, which is useful for plotting
functions. Matplotlib is also used here to export the graphs of x(t) and v(t) to a file.
Exporting output (either arrays or images) to a file is useful, so that the output of the
program can be used in other places without having to rerun the original program. After
importing the libraries, we then define variables in the next block of code. Notice that we
are choosing the step size to be dt = 0.01, and we want to integrate equations from t = 0 to
t = 30.0 (tmax ). The variable num steps = 3000 is the number of steps of size 0.01 between
0 and 30 and will be used to determine the length of arrays.
After defining variables in Algorithm 2, we defined the second term in the right-hand side
of each equation in (2.7.9) as the functions dxdt and dvdt, respectively. Although it is not
necessary to define those functions, it makes the lines in the upcoming for-loop cleaner and
easier to read. Readability is important when writing algorithms, so that others can follow
your work and so that you can follow your own work, if you need to reuse the algorithm
later.
Next, we need to define the arrays x and v, which will contain our solutions x(t) and
v(t), respectively. We need the arrays to have enough elements to cover the time range from
t = 0 to 30. We create the arrays using the command zeros from the NumPy library. The np
in front of the command zeros tells the Python interpreter that the command zeros comes
from the NumPy library. The command zeros creates an array in which all elements are 0.
The length of the array is determined by the argument of the command; in this case, the
arrays x and v have a length num steps + 1. This may seem strange, but we can explain
it using Python’s method for indexing arrays. Python begins counting using 0, so the first
element of an array has an index of 0 (note that some languages, like Mathematica, give
the first element of an array the index of 1). There are 3000 steps of size 0.01 between t = 0
and t = 30, and we need one more additional step for the time t = 0. Hence, in order to span
the time from 0 to 30 in steps of 0.01, we need 3001 elements in our arrays. By initially
assigning all elements the value of zero, we have automatically included the initial condition
as the first element of each array.
Next, we use a for-loop to perform Euler’s method. The for command tells Python to
start a for-loop. The contents of the for-loop (the two-indented lines below the for statement)
are repeated multiple times. The statement, for i in range(1,num steps+1), tells the Python
interpreter that the value of i starts at i = 1. With each iteration of the loop, the value of
i increases by 1. The loop stops when i = num steps. The range command produces a list
of values from 1 up to, but not including, num steps + 1. Therefore, the loop is iterated
num step times. In each iteration of the loop, we calculate a new element of each array x
and v. Hence, in the first iteration of the loop i = 1, and we compute x[1] and v[1], and the
value of i is changed to 2. In the next iteration, i = 2, and we find x[2] and v[2] and the
value of i is changed to 3. This process continues until i has the value num steps + 1. Note
that the value of i is increased after x[i] and v[i] are computed, so once i = num steps + 1
the loop stops without computing x[num steps +1] and v[num steps + 1].
In order to plot the functions x and v, we need to tell Python to which time each array
element corresponds. The time variable contains a list of times from 0 to 30 in steps of
0.01 using the command arange, which is similar to the for command and ends one step
before the second argument, hence the tmax + dt term. The algorithm ends by creating the
necessary plots and saving them as an encapsulated postscript file.
i i
i i
i i
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
#d e f i n e f u n c t i o n s
d e f dxdt ( v e l ) :
return vel
d e f dvdt ( v e l ) :
r e t u r n g − k/m ∗ v e l
#d e f i n e i n i t i a l c o n d i t i o n s and a r r a y s
x = np . z e r o s ( num steps +1)
v = np . z e r o s ( num steps +1)
time = np . a r a n g e ( 0 , tmax+dt , dt )
i i
i i
i i
The output of Algorithm 2 is shown in Figure 2.5. Notice that the results obtained from
the Python program look similar to the results obtained from Mathematica. That is a good
sign! When we plot both results on the same graph, we find that the graphs are actually
identical. When developing your own numerical analysis programs, such as the Euler method
algorithm, it is often wise to check the results of your new algorithm with another algorithm
that you suspect works fine. In this case, the Mathematica DSolve command produced the
same analytical form as we obtained by hand. Therefore, we verified Mathematica’s result
with our own calculation and then verified our Euler method result (found using Python)
with that of Mathematica’s result. Ultimately, we can be confident that our Euler method
program works well, at least in the case of this one particular ODE. Furthermore, our Euler
method’s success here makes us more confident of the results we would obtain by using our
algorithm on new problems.
1200 50
1000 40
800 30
v, m/s
x, m
600
20
400
200 10
0 0
0 10 20 30 0 10 20 30
time, s time, s
The ability to solve differential equations numerically allows you to address interesting
problems that you could not do by hand. The next example involves three different forces
acting on a system, it demonstrates the power of Euler’s method.
i i
i i
i i
Solution:
We will be using Euler’s method to solve the differential equation. In order to use
Euler’s method, we need to turn the second-order ODE into two coupled first-order ODEs.
This is done by using v = ẋ:
ẋ =v
(2.7.12)
v̇ = − 3x − 0.5v 2 + cos(t)
In the algorithm below, we removed the comments and graphing commands, but those
are similar to the previous problem involving linear air resistance. Furthermore, we added
the resulting graphs of x(t) and v(t) after the algorithm.
import numpy a s np
dt = 0 . 0 1
tmax = 4 0 . 0
num steps = i n t ( tmax/ dt )
d e f dxdt ( v e l ) :
return vel
d e f dvdt ( pos , v e l , t ) :
r e t u r n np . c o s ( t )−3∗ pos − 0 . 5 ∗ v e l ∗∗2
x = np . z e r o s ( num steps )
v = np . z e r o s ( num steps )
f o r i i n r a n g e ( 1 , num steps ) :
x [ i ] = x [ i −1] + dt ∗ dxdt ( v [ i −1])
v [ i ] = v [ i −1] + dt ∗ dvdt ( x [ i −1] , v [ i −1] , i ∗ dt )
1.0
2
0.5
1
velocity, m/s
position, m
0.0
−0.5 0
−1.0
−1
−1.5
−2.0 −2
0 10 20 30 40 0 10 20 30 40
time, s time, s
Notice that we couldn’t pass the loop variable i into the function computing the cosine.
By itself, i is just an integer and does not represent time. The time at step i is idt.
i i
i i
i i
The position versus time graph is irregular. Notice that the particle position is not
continuously increasing or decreasing, nor is it periodic (it does not exactly repeat itself).
The differential equation governing this particle’s behavior is nonlinear, and nonlinear
systems can display interesting behavior. We’ll discuss nonlinear systems in more detail
in Chapter 13.
Example 2.8 can also be used to demonstrate the challenge of choosing a good value
for dt. By choosing dt = 0.01 in Example 2.8, we obtained a good solution for 0 ≤ t ≤ 40.
However, nonlinear equations can be tricky. If we choose dt = 0.02, the Euler method will
become unstable and the solution goes to infinity (try it!).
d2 x dv dv
F =m =m = mv
dt2 dt dx
F = m dv
dt is a first-order ODE, because the highest derivative is a first derivative.
2
F = m ddt2x is a second-order ODE, because the highest derivative is a second derivative.
The number of arbitrary constants in the general solution of an ODE is equal to the
order of the ODE.
The magnitude of the air resistance f (v) has, in general, both a linear and a quadratic
component in v.
f (v) = bv + cv 2
We find the terminal velocity for a falling object by setting the air resistance force to
be equal to the weight of the object, f (v) = bv + cv 2 = mg.
i i
i i
i i
When the force is a function of the position only, F (x), we can integrate Newton’s
second law in the form:
dv
F (x) = ma = mv
dx
to obtain the function v(x). Then integrate once more the v(x) = dx/dt equation to obtain
the position x(t).
An important system in mechanics is the mass-spring system described by Hooke’s
Law:
d2 x
F (x) = m 2 = −kx
dt
With the initial conditions v(0) = 0 and x(0) = x0 at t = 0, this equation has solutions
x(t) and v(t) given by:
r ! r !
k k
x(t) = x0 cos t v(t) = dx/dt = −ωx0 sin t
m m
and the mass oscillates about the equilibrium position x0 with the natural frequency ω of
the oscillator r
k
ω=
m
Euler’s method is a simple numerical technique for integrating differential equations
which uses the derivative of the function to approximate future values of the differential
equation’s solution.
d2 x
= −ω 2 x (2.9.2)
dt2
where x(t) is the position of a particle as a function of time. One method of solving
second-order differential equations is by a using trial solution. In this case, the equation
says that the second derivative of the function x(t) returns itself with a constant and
negative sign. What functions behave this way? Certainly sine and cosine, however an
exponential solution would also work, x(t) = Aeλt . Insert into the differential equation
the exponential function as a trial solution. What values of λ are required in order for
Aeλt to be a solution to the differential equation? What is the general solution to the
differential equation?
i i
i i
i i
4. Prove that the result from Problem 3 holds for any linear nonhomogeneous ODE.
10. Consider the double Atwood machine shown in Figure 2.6. Compute the acceleration
of the mass m1 in terms of g, M , and m2 . Ignore friction and assume the pulleys are
massless.
i i
i i
i i
m
m 2
1
Figure 2.6: Problem 2.10: A double Atwood machine.
15. Find the terminal velocity of an object experiencing a linear and quadratic air resis-
tance, as well as a gravitational force, i.e., F (v) = βv + γv 2 − mg.
i i
i i
i i
16. An object experiences a force, F (v) = −cv 2 . What are the units of c if the velocity of
the particle is measured in meters per second and its mass is measured in kilograms?
If the particle starts at the origin with an initial velocity, v(0) = v0 , find the particle’s
velocity and position as a function of time. Use a computer algebra system to verify
your results.
17. A ball is thrown upwards with an initial velocity of 10 m/s. The ball has a mass
m =0.10 kg and experiences a drag force F = −bv that scales linearly with the ball’s
velocity such that the magnitude of b = 0.2:
i i
i i
i i
(b) How long will it take her to fall if the parachute opens? Assume that the parachute
opens immediately and that the drag force on Sarah with the parachute open is
described by,
1
f = − cd ρAv 2 (2.9.3)
2
where the drag coefficient, cd = 1.5 for the parachute, ρ = 1.22 kg/m3 is the
density of air, and A = 46 m2 is the area of the parachute.
25. Consider a falling object whose mass is 1 kg and experiences an air resistance force
of F (v) = −c |v| v, where c = 0.5 Ns2 /m2 . Assuming the falling object starts at rest,
use a computer to calculate the time it takes the object to fall a distance of 1 km.
Does the velocity of the falling object ever become constant? In other words, is there
a terminal velocity for this force?
Section 2.6: Position-Dependent Forces
26. A particle has a velocity v(x) = ax−n . Find the force acting on the particle. If the
particle starts at rest at the origin, find its position and velocity as functions of time.
27. Find the velocity v(x) and position x(t) of a particle starting at rest at the origin
experiencing the force F = a + bx with a, b > 0.
28. A particle experiences a force causing it to move in the +x-direction with velocity,
v(x) = cxe−bx . Find the work done by the force if the particle starts at the origin and
travels to infinity (note: in this case, infinity is a useful mathematical tool meaning
“very far”). The work is defined by:
Z x
W= F · dx
0
29. Compute the value of the acceleration due to gravity as a function of an object’s
height y above the Earth’s surface. Plot your result from y = 0 to 25 km. Hint: Use
Newton’s universal law of gravitation:
Gm1 m2
Fg = − (2.9.4)
r2
where G = 1.67 × 10−11 Nm2 /kg2 is the universal gravitation constant, m1 and m2 are
the masses of the interacting bodies, and r is the distance between the centers of m1
and m2 . Assume that the falling object is a point particle and ignore air resistance.
30. As found in Problem 29, the acceleration of a falling object is not constant with height.
It turns out that the acceleration varies as:
9.8 m/s2
g= 2 (2.9.5)
1 + Rye
where Re = 6.37 × 106 m is the radius of the Earth, y is the altitude of the object
above the Earth’s surface, and g is in m/s2 . Compute the time it takes for a m =1
kg object with a diameter D =1 m to fall from an altitude of 25 km and land on the
Earth’s surface, assuming quadratic air resistance. What is the object’s final velocity?
You will need a computer to solve this problem.
i i
i i
i i
31. A particle experiences a force F = cvx, where c > 0. If at t = 0, the particle is passing
through the origin with a velocity v = v0 , find the particle’s position as a function of
time.
i i
i i
i i
CHAPTER 3
In the previous chapter, the kinematics and dynamics aspects of particle motion were exam-
ined in one dimension by using the scalar position parameter x. In this chapter, we will first
describe the motion of particles in two and three dimensions using Cartesian coordinates
(x, y, z). Describing motion in two and three dimensions requires the use of vectors. After
reviewing the basic properties of vectors and introducing the dot and cross products, we
will introduce vector derivatives. Vector derivatives are necessary in order to compute the
velocity and acceleration vectors. We then describe polar, cylindrical, and spherical coordi-
nates, and demonstrate how to describe a particle’s position, velocity, and acceleration in
those coordinate systems. Finally, we conclude the chapter with special vector derivatives
commonly used in all fields of physics, such as the gradient, the divergence, and the curl.
The magnitude of a vector is often thought of as its length and would correspond to a
quantity measured from the vector. For example, suppose A represents the position of
61
i i
i i
i i
a particle located at the point A = 3 mî + 4 mĵ. The magnitude of A is |A| = 5 m and
represents the distance between the particle and the origin.
It is often useful to know the direction in which a vector points. This is easily found
for simple vectors like A = 3î, which points along the x-axis. However, other vectors like
A = aî + bĵ do not lie along either the x- or y-axes. In that case, we need a unit vector which
points in the same direction as the vector A but has a length of one. We can obtain a unit
vector  in the direction of any vector A, by dividing A by its magnitude:
A A
 = =q (3.1.3)
|A| A2x + A2y + A2z
If we know the coordinates of two points P1 and P2 in space, then we can find the vector
connecting these two points by subtracting their respective x-, y-, z-components. For exam-
ple, in Figure 3.1, let r1 be the vector from the origin to P1 , and r2 the vector from the
origin to P2 . The vector r12 from P1 to P2 , is given by (see Figure 3.1):
r12 = r2 − r1 (3.1.4)
Now that we have covered some of the basic concepts behind vectors, we can apply
them to the basic descriptors of motion. As discussed in Chapter 1, the position vector
function r(t) = (x(t), y(t), z(t)) for a particle in three dimensions as described by Cartesian
coordinates is written as:
Position Vector
This is a vector-valued function, and x(t), y(t), z(t) are the scalar functions representing
the position along the x-, y-, and z-axes at time t, respectively. The vector r(t) has its tail
at the origin, and its head at the position of the particle at time t. The magnitude of the
position vector is equal to the distance between the particle and the origin.
i i
i i
i i
Figure 3.1: (a) The Cartesian components of vector A = Ax î + Ax ĵ + Ax k̂; (b) The vector
r12 = r2 − r1 connecting two points P1 and P2 in space.
The derivative of the position vector is the velocity v(t) of the particle. To find the
velocity, we need to know how to differentiate vectors. Figure 3.2 illustrates the derivative
of a vector function r. Without loss of generality, r can be thought of as a position vector,
however, what follows is true for any vector. For the derivative of a vector function to
exist, the vector function must be a continuous function of the scalar variable t. The vector
function r(t) is shown as the dashed line in Figure 3.2, where the vector r is shown at times t
and t + ∆t. The vector ∆r = r(t + ∆t) − r(t) represents the change in the value of the vector
function and will decrease in magnitude as ∆t → 0. Therefore, we can find the derivative
using the standard formula:
dr(t) d
= x(t)î + y(t)ĵ + z(t)k̂ (3.1.9)
dt dt
d d d
= x(t)î + y(t)ĵ + z(t)k̂ (3.1.10)
dt dt dt
However, the unit vectors in Cartesian coordinates are constant vectors. Therefore, the
velocity vector can be found as:
i i
i i
i i
Velocity Vector
dr dx dy dz
v(t) = = î + ĵ + k̂ (3.1.11)
dt dt dt dt
Because the velocity vector is the derivative of the position vector, the velocity vector
is going to be always tangent to the path followed by the particle.
Figure 3.2: The vector function r is illustrated here at times t and t + ∆t. In the limit
∆t → 0, the vector ∆r/∆t becomes the derivative vector function dr/dt, which is tangent
to the curve.
The magnitude v(t) of the velocity vector v(t) is called the speed:
s
2 2 2
dx dy dz
v(t) = |v(t)| = + + (3.1.12)
dt dt dt
Equation (3.1.12) can also be written in terms of the arc length s along the curve C describ-
ing the motion, as measured from some initial position along the curve. In order to use the
arc length, we need to define an infinitesimal displacement along the curve,
ds = dxî + dy ĵ + dz k̂ (3.1.13)
The differential distance ds is the magnitude of the infinitesimal displacement and is given
by q
ds = (dx)2 + (dy)2 + (dz)2 (3.1.14)
Equation (3.1.12) can now be written in terms of the arc length s, as:
ds
v(t) = |v(t)| = (3.1.15)
dt
A unit tangent vector T along the curve C describing the motion is obtained by dividing
the velocity vector v(t) by its magnitude v:
v(t)
T= (3.1.16)
v(t)
i i
i i
i i
Note that the unit tangent vector T = v̂ points along the direction of the velocity and is
always tangent to the curve C. An example of how to calculate the unit tangent vector is
given in Example 3.2.
Similarly, the derivative of the velocity vector dv/dt is the acceleration vector function,
which in Cartesian coordinates is:
Acceleration Vector
dv d2 x dy 2 d2 z
a(t) = = 2 î + 2 ĵ + 2 k̂ (3.1.17)
dt dt dt dt
q
v(1) = [2.0 cos(2.0)]2 + (2)2 = 2.2 m/s
q
a(1) = [4.0 sin(2.0)]2 + (2)2 = 4.2 m/s2
(b) The unit vector T, which is tangent to the curve C at any time t, is found from:
1
v(t) dx(t) dy(t) dz(t)
T= = î + ĵ + k̂
v(t) v(t) dt dt dt
By evaluating the derivatives and substituting the speed v(t) from part (a):
1 1
T= [2 cos(2t)] î + 2tĵ = q [2 cos(2t)] î + 2tĵ
v(t)
[2 cos(2t)]2 + (2t)2
The next example shows how to calculate the position vector when one is given an
acceleration vector. As you might expect, the problem involves integrating the acceleration.
Example 3.3: Finding position and velocity from acceleration vector.
The acceleration of a particle moving in two dimensional space is given by the vector:
i i
i i
i i
where the vector a is measured in m/s2 . Find the velocity vector v(t) and the position
vector r(t). Assume that the particle is stationary and is at the origin at time t = 0.
Solution:
Since a(t) = dv(t)/dt, we can find the velocity and position by integration. Just like dif-
ferentiating a vector, integrating a vector can be done one component at a time:
Z
v(t) = a(t)dt
Z
= sin (2t) î + cos (2t) ĵ dt
Z Z
= sin (2t) îdt + cos (2t) ĵdt
which gives d = − 41 m ĵ. Again, we have used the fact that the unit vectors î and ĵ in
Cartesian coordinates are constant, and therefore can be placed outside of each integral.
The three vector functions r(t), v(t), and a(t) depend on the single parameter t, rep-
resenting time. If we are given the explicit functions x(t), y(t), z(t) then the path of the
moving particle can be plotted in a so-called parametric plot in three dimensions using a
computer, see Examples 3.4 and 3.5 that follow. A parametric plot uses the functions x(t),
y(t), and z(t) to graph a curve that represents the path taken by a particle. Each point in
the curve corresponds to the x-, y-, and z-coordinates for various times t.
In the case of two-dimensional motion, if we are given explicit functions x(t) and y(t), one
may be able to eliminate the variable t from the simultaneous equations x = x(t), y = y(t).
If one of these equations can be solved for t, the expression obtained can be substituted
into the other equation to obtain an expression for y(x). In some cases, there is no single
equation in closed form that is equivalent to the parametric equations. A similar process
can be used in three dimensions in order to get z(x, y). Examples 3.4 and 3.5 show how to
create parametric plots in two and three dimensions in Mathematica and Python.
i i
i i
i i
y
1.2
0.6
Out[ ]=
x
-2.4 -1.2 1.2 2.4
-0.6
-1.2
i i
i i
i i
This is the equation of an ellipse on the xy-plane with center at (d, e). Since z = ct, one
then expects the motion to be in the form of an elliptical helix moving along the z-axis
with a uniform speed c.
Next we present Python and Mathematica codes for plotting the motion, with the
numerical values a = 2, b = 4, c = 1, d = 1, e = 2, T = 12. The Python code is below.
import m a t p l o t l i b a s mpl
from m p l t o o l k i t s . mplot3d import Axes3D
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
mpl . rcParams [ ’ l e g e n d . f o n t s i z e ’ ] = 10
p l t . show ( )
i i
i i
i i
A · B = Ax Bx + Ay By + Az Bz (3.2.1)
A · B = AB cos θ (3.2.2)
where A and B are the magnitude of the vectors A and B, respectively, and θ is the angle
between the two vectors as shown in Figure 3.3. Notice that the dot product is multiplying
the magnitude of B by the component of A, which is parallel to B. A physical example
involving the dot product is the definition of work. The work done by a constant force can be
computed as W = F · r, where a force F is applied to a particle along a displacement r. The
component of the force doing work on the particle is parallel to the particle’s displacement.
The dot product is used to compute the work because we need the component of F parallel
to r.
Equation (3.2.2) can be used to obtain the angle between two vectors:
A·B Ax Bx + Ay By + Az Bz
cos θ = = q q . (3.2.3)
AB A2x + A2y + A2z Bx2 + By2 + Bz2
The dot product can also be interpreted geometrically as the scalar projection of vector
A on the direction of vector B, as shown in Figure 3.3. This is useful if you want to know
the component of the vector A along the direction of the vector B. First, we create the
unit vector B̂ =B/ |B| pointing in the direction of B. The scalar projection is equal to
AB = A · B̂ = |A| cos θ. The vector projection of A onto B is, projB A = AB B̂.
Figure 3.3: The dot product represents the product of the scalar projection of vector A on
the unit vector B̂ in the direction of vector B: A · B̂ = |A| cos θ.
i i
i i
i i
There are several algebraic relationships for the dot products of two vectors A and B,
and most of them are intuitive. The following properties hold if A, B, and C are real vectors
and c is a scalar.
The dot product is commutative:
A · (B + C) = A · B + A · C (3.2.5)
A · (cB) = c (A · B) (3.2.6)
The geometric interpretation of the orthogonality of vectors is that they are perpendicular,
i.e., θ = π/2. Finally, the product rule applies to the dot product:
d dA dB
(A · B) = ·B+A· (3.2.8)
dt dt dt
OUTPUT: 0
The Mathematica code that follows uses the command Dot[A, B], and the vectors
A, B are entered in the form of lists of their Cartesian components.
i i
i i
i i
√
At t = 1 s we find the unit tangent vector T = 2î + ĵ / 5.
(b) Since T is a unit vector with a magnitude of 1, the dot product |T|2 = T · T = 1.
By differentiating |T|2 with respect to time t, we obtain:
dT dT dT
T + T = 2T =0
dt dt dt
Therefore the vector dT/dt is perpendicular to the tangent unit vector T at any time
t, and is found from: " #
dT d 2tî + ĵ (2) î + (−4t) ĵ
= √ =
dt dt 4t + 1
2
(4t2 + 1)
3/2
(c) From part (b) we obtain at t = 1 s: dT dt = 5
−3/2 2î − 4ĵ , and its magnitude is:
vu !2 !2 √
dT u 2 −4 20
= + =
t
dt (5)3/2 (5)3/2 (5)3/2
The desired unit vector N in this perpendicular direction is obtained by dividing dT/dt
by its magnitude:
2î−4ĵ
dT
(5)3/2 2î − 4ĵ
N= dt
= √ = √
| dT
dt |
20 20
(5)3/2
i i
i i
i i
a vector C that is perpendicular to both vectors A and B as shown in Figure 3.4(a). The
direction of the cross product vector A × B is given by the right-hand rule shown in Figure
3.4(b). The magnitude of this vector is equal to the area of the parallelogram formed by
vectors A and B, as shown in Figure 3.4(c).
Figure 3.4: (a) The cross product in respect to a right-handed coordinate system; (b) Finding
the direction of the cross product by the right-hand rule; (c) The area of a parallelogram
formed by two vectors A and B, is equal to the magnitude of the cross product of these
two vectors.
a × b = ab sin θ n̂ (3.2.9)
where θ is the angle between a and b (0 ≤ θ ≤ π), a = |a| and b = |b| are the magnitudes of
vectors a and b, and n̂ is a unit vector perpendicular to the plane containing a and b in the
direction given by the right-hand rule as in Figure 3.4(a). If the vectors a and b are parallel
(i.e., the angle θ between them is either 0 radians or π radians), then the cross product of
a and b is the zero vector 0.
The magnitude of the cross product is the product of the magnitude of the vector a,
and the component of b perpendicular to a. You may recall from your introductory physics
course that you can compute the torque applied to a door by using the formula N = r × F,
where N is the torque applied by the force F, and r is the vector which points from the
axis of rotation to the point where the force is applied. For a door, the vector r lies along
the door, pointing from the hinges to the edge of the door, where the force to open it is
applied. We know that the component of the force responsible for the rotation of the door
is perpendicular to the door, and that the component of the force parallel to the door is not
causing the rotation. Therefore, the cross product is used to compute the torque, because
we need the component of the force perpendicular to r.
The cross product of two vectors a = [a1 , a2 , a3 ] and b = [b1 , b2 , b3 ] can also be repre-
sented by the determinant of a formal matrix:
î ĵ k̂
a × b = a1 a2 a3 (3.2.10)
b1 b2 b3
i i
i i
i i
This determinant can be expanded by using cofactor expansion along the first row to yield:
a a3 a1 a3
a × b = 2 ĵ + a1 a2 k̂ (3.2.11)
î −
b2 b3 b1 b3 b1 b2
which gives the components of the resulting cross product vector directly:
Like the dot product, the cross product has some familiar algebraic identities. The cross
product is distributive with addition:
a × (b + c) = a × b + a × c (3.2.13)
This property is easily concluded by considering the right-hand rule; when the two vectors
are interchanged, the sign of the cross-product vector is reversed. The product rule for
derivatives also applies to the cross product:
d da db
(a × b) = ×b+a× (3.2.16)
dx dx dx
The following are the Python and Mathematica codes for the cross product.
OUTPUT:
− 5∗N. x + 2∗N. y − 7∗N. z
i i
i i
i i
To find the unit vector, divide this vector by its magnitude. The magnitude of a × b
is: q √
|a × b| = 12 + 52 + (−14)2 = 222
The desired unit vector is
In the Mathematica code below we use the Cross command to find the cross product,
and we use the Norm[c] command to find the magnitude of the given vector.
a = {1, 3, −1};
b = {4, −2, 1};
c = Cross[a, b]
OUTPUT: {1, −5, −14}
Norm[c] √
OUTPUT: 222
unitc = c/Norm[c]
n q o
OUTPUT: √ 1 , − √ 5 , −7
222 222
2
111
A useful mathematical construction involving the cross product is the triple product,
which is used in two different forms.
The scalar triple product of three vectors is defined as
a · (b × c) (3.2.17)
Notice that this expression means that we first calculate the cross product of vectors b
and c, followed by the dot product between the vectors a and b × c. The magnitude of the
triple scalar product represents the volume V of the parallelepiped with edges a, b, and
c. Since these three vectors a, b, c can be interchanged, the following three scalar triple
products are equal to each other and to the volume V :
a · (b × c) = b · (c × a) = c · (a × b) = V (3.2.18)
The vector triple product is the cross product of a vector with the result of another cross
product and is related to the dot product by the following identity:
a × (b × c) = b(a · c) − c(a · b) (3.2.19)
i i
i i
i i
The mnemonic “BAC minus CAB” is often used to remember the order of the vectors
in the right-hand side of this equation. This formula is used in physics to simplify vector
calculations. Another useful identity relates the cross product to the scalar triple product:
(a × b) × (a × c) = (a · (b × c))a (3.2.20)
OUTPUT: True
OUTPUT: True
i i
i i
i i
r̂
P
rsin r
rcos X
Figure 3.5: The radius r and angle θ locate a particle in polar coordinates, and they also
represent a complex number: z = x + iy = r(cos θ + i sin θ) = r eiθ . The two unit vectors r̂
and θ̂θ (in blue in the e-book) in polar coordinates change direction as the particle moves in
the xy plane along a curved path (in red in the e-book).
The relationship between Cartesian (x,y) coordinates and the polar coordinates r, θ is:
x = r cos θ (3.3.1)
y = r sin θ (3.3.2)
r = x2 + y 2 (3.3.3)
p
We can also write the position vector as a complex number in the following form by
using Euler’s equations:
By adding and subtracting the last two equations, we obtain Euler’s formulas for sine
and cosine functions. These equations provide a powerful connection between analysis and
trigonometry and should be memorized.
Euler’s Identities
eiθ + e−iθ
cos θ = (3.3.7)
2
i i
i i
i i
eiθ − e−iθ
sin θ = (3.3.8)
2i
eiθ = cos θ + i sin θ (3.3.9)
e−iθ = cos θ − i sin θ (3.3.10)
While polar coordinates can be convenient for describing the position of a particle,
especially a particle constrained to move along an arc or a circle (red in the e-book in
Figure 3.5), the unit vectors require special attention. In Cartesian coordinates, the unit
vectors î and ĵ always point along the x- and y-axes respectively, regardless of the particle’s
location. However, the directions of the unit polar vectors r̂ and θ̂ change with the location
of the particle. This can be clearly seen in Figure 3.5.
The unit vector r̂ always points in the radially outward direction. The position vector
of the particle moving in the xy-plane can be written in terms of the unit vector r̂ as:
r = rr̂ (3.3.11)
In order to find the velocity and acceleration in polar coordinates, we first need to write
the unit vectors r̂ and θ̂θ in terms of î and ĵ. Using the geometry shown in Figure 3.5, we
can write:
i.e., θ̂·r̂ = 0.
The unit vectors r̂ and θ̂ vary in direction with time, and we can calculate their time
derivatives by using the chain rule:
d d h i dθ dθ dθ dθ
r̂ = cos θî + sin θĵ = − sin θ î + cos θ ĵ = − sin θî + cos θ ĵ (3.3.14)
dt dt dt dt dt dt
By working in a similar manner we obtain the time derivative of the unit vector θ̂:
dθ̂θ dθ
= − r̂ (3.3.16)
dt dt
We can now calculate the velocity vector in polar coordinates by taking the first time
derivative of the position vector r:
dr d dr dr̂
v= = (rr̂) = r̂ + r (3.3.17)
dt dt dt dt
By substituting (3.3.15) into this equation, we obtain the velocity vector in polar coor-
dinates:
dr dr dθ
v= = r̂ + r θ̂θ (3.3.18)
dt dt dt
i i
i i
i i
The two components of the velocity in this equation have distinct physical meanings.
dt r̂ represents the radial component of the velocity, which points along
The first term vr = dr
the radius r. The second term v θ = r dθ dt θ̂ = rω θ̂ is proportional to the angular velocity
θ θ
ω = dθ/dt, and represents the velocity of the particle along the θ̂θ direction. This second
component represents the velocity of uniform circular motion.
From (3.3.18), we obtain the infinitesimal displacement ds between two points in space
in polar coordinates by canceling the differential time dt from both sides of this equation:
Based on this expression, the infinitesimal area element dA in polar coordinates can then
be written as:
dA = dxdy = rdθdr (3.3.20)
and we can calculate double integrals of any function f (x, y) in polar coordinates by using
the transformation:
Z Z Z Z
f (x, y) dxdy = f (r cos θ, r sin θ) r dr dθ (3.3.21)
In summary, the unit vectors in polar coordinates and their derivatives are given by:
The following example shows how to create polar plots in Python and in Mathematica.
i i
i i
i i
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
r = np . a r a n g e ( 0 , 2 , 0 . 0 1 )
t h e t a = 2 ∗ np . p i ∗ r
ax = p l t . s u b p l o t ( 1 1 1 , p r o j e c t i o n= ’ p o l a r ’ )
ax . p l o t ( t h e t a , r )
ax . s e t r m a x ( 2 )
ax . s e t r t i c k s ( [ 0 . 5 , 1 , 1 . 5 , 2 ] )
ax . g r i d ( True )
ax . s e t t i t l e ( ” Example o f a p o l a r p l o t i n Python ” , va= ’ bottom ’ )
p l t . s a v e f i g ( ’ p o l a r p l o t e x a m p l e . png ’ , d p i =600)
p l t . show ( )
Polar coordinates are very useful in describing the motion of celestial bodies around
each other, as we will explore in detail in a later chapter. The following example shows
how polar coordinates can simplify the mathematical description of the elliptical motion of
planets around the Sun.
Example 3.12: Ellipses in polar coordinates.
The equation of an ellipse with semi-major axes a, b in Cartesian coordinates is:
(x/a)2 + (y/b)2 = 1
or in parametric form
x = a cos θ y = b sin θ.
i i
i i
i i
(a) Use the symbolic capabilities of Mathematica to show that the general equation for
this ellipse in polar coordinates is:
ab
r= √ (3.3.25)
a − c2 cos2 θ
2
√
where c = a2 − b2 (with a > b). √
(b) Show that by shifting the x-axis by the quantity c = a2 − b2 , the equation of the
ellipse in polar coordinates simplifies to:
a 1 − e2
b2 /a
r= = (3.3.26)
1 + (c/a) cos θ 1 + e cos θ
(c) Plot the above expressions in Mathematica using the PolarPlot command.
Solution:
(a) The following Mathematica code performs the necessary algebraic steps for deriving
(3.3.25), and then plots the ellipse with parameters
√ a = 2, b = 1. After defining the para-
metric forms x = a cos θ and y = b sin θ = b 1 − cos2 θ, the code uses the Solve command
to find the distance r as a function of cos θ. The command “replace with”/. is used to
define a parameter rellipse for the desired equation r = r(θ). Since there are two solutions
obtained from the Solve command, we chose the positive solution by selecting the second
(positive) square root.
x = r ∗ Cos[θ]; y = r ∗ Sqrt[1 − Cos[θ]∧ 2];
sol = Solve[(x/a)∧ 2 + (y/b)∧ 2 == 1, r]
OUTPUT: r → −q 1
, r → q 1
1 + Cos[θ]2 − Cos[θ]2 1 + Cos[θ]2 − Cos[θ]2
b2 a2 b2 b2 a2 b2
rellipse = r/.sol[[2]]
OUTPUT: q 1
1 + Cos[θ]2 − Cos[θ]2
b2 a2 b2
0.8
0.4
In[ ]:=
x
-1.8 -1.2 -0.6 0.6 1.2 1.8
-0.4
-0.8
(b) The following code carries out the algebra for deriving (3.3.26), and plots the ellipse.
As the polar plot shows, the resulting ellipse from (3.3.26) is shifted along the x-axis
i i
i i
i i
b2
OUTPUT: √
a+ (a−b)(a+b)Cos[θ]
1.0
0.5
x
-4 -3 -2 -1 1
-0.5
-1.0
The
h 2two terms2in
i this equation again have distinct physical meanings. The first term
ar = dt2 − r dt
d r dθ
r̂ represents the radial component of the acceleration, which points
h 2 i
along the radius r. The first part of this radial term is ddt2r r̂ and represents the rate of
2
change of the radial velocity dr/dt. The second part of this radial term is −r dθ
dt = −rω 2 ,
d2 r
and points in the opposite direction from the dt2
term, and represents the centripetal
d2 θ
acceleration ac = v 2 /r. The second term aθ = (r dt2 + 2 dr
dt dt )θ̂ also contains two terms; the
dθ θ
2
first term r ddt2θ = r dω
dt = rα is the angular acceleration along the θ̂ direction. The last term
θ
represents the Coriolis force and will be discussed in Chapter 10.
In summary, the position, velocity, and acceleration vectors in polar coordinates are given
by:
r = rr̂ (3.3.28)
dr dr dθ
v= = r̂ + r θ̂θ (3.3.29)
dt dt dt
i i
i i
i i
Figure 3.6: Cylindrical coordinates (r, θ, z) and the corresponding unit vectors r̂, θ̂, ẑ.
The transformation from Cartesian to polar coordinates is made by the following rela-
tions:
From Cylindrical to Cartesian Coordinates
x = r cos θ (3.3.31)
y = r sin θ (3.3.32)
z=z (3.3.33)
We can define three unit vectors r̂, θ̂θ , ẑ in cylindrical coordinates, which are shown in
Figure 3.6. Since cylindrical coordinates are essentially an extension of polar coordinates in
three dimensions, we can write the unit vectors immediately by analogy with the equations
in the previous section.
i i
i i
i i
Similar to the unit vectors in polar coordinates, the unit vectors in cylindrical coordinates
change direction with the location of the particle, with the exception of ẑ, which is identical
to k̂ and does not change direction. Note also that the cylindrical unit vectors are related
to the Cartesian unit vectors by a rotation matrix of the form:
cos θ sin θ 0
r̂ î
θ̂θ = − sin θ cos θ 0 ĵ (3.3.40)
ẑ 0 0 1 k̂
We can compute the position and velocity vectors in cylindrical coordinates in a manner
similar to that used for polar coordinates. The difference is that for cylindrical coordinates
there is an additional direction, ẑ. However, ẑ does not change direction with position and
is therefore constant. Hence, the position, velocity, and acceleration vectors in cylindrical
coordinates are given by:
r = rr̂ (3.3.41)
dr dθ dz
v= r̂ + r θ̂θ + ẑ (3.3.42)
dt dt dt
" 2 #
dv 2
d r dθ 2
d θ dr dθ d2 z
a= = − r r̂ + (r + 2 ) θ
θ̂ + ẑ (3.3.43)
dt dt2 dt dt2 dt dt dt2
This section concludes with some additional useful expressions. By analogy to (3.3.19),
which was derived for polar coordinates, the infinitesimal displacement ds between two
points in space in cylindrical coordinates can be written as:
When computing infinitesimal areas, it is important to keep in mind the specific area on
the cylinder that is relevant to your problem. It is common to work with the infinitesimal
area element that is on the surface of the body of the cylinder. The infinitesimal area element
dA on the surface of a cylinder of radius r, can be written as:
dA = rdθdz (3.3.45)
Keep in mind, however, that if one is interested in the top or bottom of the cylinder, then
the polar infinitesimal area element, dA = rdrdθ, applies. An infinitesimal volume element
in cylindrical coordinates is given by:
i i
i i
i i
and we can calculate triple integrals of any function f (x, y, z) by using the transformation
into cylindrical coordinates:
ZZZ ZZZ
f (x, y, z) dxdydz = f (r cos θ, r sin θ, z) r dr dθ dz (3.3.47)
These triple integrals will become useful in calculating, for example, the center of mass and
moment of inertia of an object in cylindrical coordinates.
Figure 3.7: Spherical coordinates (r, θ, φ). The angle φ is defined on the xy-plane and can
vary between 0 to 2π, while the azimuthal angle θ can vary between 0 to π only.
i i
i i
i i
z z
cos θ = =p (3.3.53)
r x + y2 + z2
2
Similar to the unit vectors in cylindrical coordinates, the orientation of the spherical
coordinate unit vectors r̂, θ̂θ , φ̂
φ depends on the position of the particle. We define these three
unit vectors as shown in Figure 3.7.
These unit vectors r̂, θ̂θ , φ̂φ for spherical coordinates are related to the Cartesian unit
vectors by the following equations:
Unit Vectors in Spherical Coordinates
Like the cylindrical coordinate unit vectors, the spherical unit vectors are related to the
Cartesian unit vectors by an orthogonal rotation matrix:
i i
i i
i i
r = rr̂ (3.3.66)
df
f (x + dx) = f (x) + dx + · · · (3.4.1)
dx
We can ignore terms of dx2 and higher, because dx is very small. Next, we define the
differential df = f (x + dx) − f (x) and interpret it as the change of the value of f when we
change x by a small amount dx. Then from (3.4.1) we find:
df
df = dx (3.4.2)
dx
Equation (3.4.2) can be interpreted as follows. The amount that f changes when x is changed
(by a small amount), is equal to the product of the rate of change of f with respect to x,
and the amount by which x is changed (dx). This is similar to the familiar equation d = vt,
which says the distance traveled by a particle is the product of how fast it is traveling (i.e.,
the rate v at which the distance changes with time), and of the length of time t traveled.
Next, we can extend (3.4.2) to scalar functions in three dimensions. Let us consider a
scalar function f (x, y, z) in Cartesian coordinates. We would like to know how f (x, y, z)
changes when we change each of its variables by a small amount. A multivariable Taylor
series expansion gives:
∂f ∂f ∂f
f (x + dx, y + dy, z + dz) = f (x, y, z) + dx + dy + dz + · · · (3.4.3)
∂x ∂y ∂z
In this expression, ∂f ∂f ∂f
∂x , ∂y , ∂z are, of course, the partial derivatives of f with respect to
(x, y, z) respectively. Similar to the one-dimensional case, we define the differential as df =
f (x + dx, y + dy, z + dz) − f (x, y, z) and interpret it as the change of f as we change all of
its variables. Using (3.4.3), we can write in dot product form:
i i
i i
i i
∂f ∂f ∂f
df = dx + dy + dz (3.4.4)
∂x ∂y ∂z
∂f ∂f ∂f
= î + ĵ + k̂ · dxî + dy ĵ + dz k̂ (3.4.5)
∂x ∂y ∂z
Again, we get an equation similar to d = vt, where the first term is similar to a rate of
change, and the second term is a change in the coordinates. The first term in (3.4.5) can
be rewritten as:
∂ ∂ ∂
î + ĵ + k̂ f (3.4.6)
∂x ∂y ∂z
and it describes the rate of change of f in each direction. The operator appearing in (3.4.6)
is very important in physics and is called the del operator denoted by ∇:
The symbol ∇ is called nabla. The del operator is a vector operator and in (3.4.7), it is
written in Cartesian coordinates. As you might expect, the del operator takes different
forms in cylindrical and spherical coordinates. We will present the del operator in other
coordinates later in this section, but for now we will consider only Cartesian coordinates.
When the del operator is applied to a scalar function, ∇f is called the gradient of f and
is a vector function.
The Gradient of a Scalar Function f (x, y, z)
∂f ∂f ∂f
∇f = î + ĵ + k̂ (3.4.8)
∂x ∂y ∂z
df = ∇f · dr (3.4.9)
To better understand the gradient of a scalar field, we will introduce the concept of level
sets of the function f . The level set of a scalar function f is a curve in two dimensions, or
a surface in three dimensions, on which the value of f is constant. Hence, along the curve
of a level set we would have df = 0. If dr is a displacement along a level set of f , then:
df = ∇f · dr = 0 (3.4.10)
and therefore the two vectors ∇f and dr are perpendicular to each other, i.e., ∇f ⊥ dr. We
conclude that the gradient of f is perpendicular to the level sets of f .
In order to get a complete picture of the gradient, we next need to discuss the concept
of the directional derivative.
The directional derivative df /dn of a scalar field f gives the rate of change of f (x, y, z)
in the direction of the unit vector n̂ = nx î + ny ĵ + nz k̂ , and is defined as the dot product:
df ∂f ∂f ∂f
= n̂ · ∇f = nx + ny + nz (3.4.11)
dn ∂x ∂y ∂z
i i
i i
i i
Now we ask the question, in what direction does the scalar field f change most rapidly?
We can compute the dot product as:
df
= n̂ · ∇f = |n̂| |∇f | cos θ (3.4.12)
dn
where θ is the angle between n̂ and ∇f . The direction in which the directional derivative is
greatest would occur when n̂ · ∇f is maximal, i.e., when the angle θ equals zero. Therefore,
∇f points in the direction of the most rapid increase of f . Likewise, the direction of greatest
decrease in the scalar field is −∇f . As an example, consider a surface whose height above
sea level at a point (x, y) is H(x, y). The gradient of H at a point (x, y) is a vector pointing
in the direction of the steepest slope, or steepest grade at that point. The steepness of the
slope at that point is given by the magnitude of the gradient vector.
(a) Compute and plot the gradient vector of f . (b) Plot the level sets of this function
by using a contour plot to show curves of constant f .
Solution:
The gradient of f can be found by computing:
∂ ∂
e−(x +y )
2 2
∇f = î + ĵ
∂x ∂y
−(x2 +y 2 )
î − 2ye−(x +y ) ĵ
2 2
= − 2xe
Next, we can produce the vector plot and the contour plot in Mathematica. The code
needed to compute ∇f , the contour plot, and the vector field plot is shown below.
f = Exp − x2 + y 2 ;
i i
i i
i i
1.0
0.5
0.0
-0.5
-1.0
-1.0 -0.5 0.0 0.5 1.0
The contour curves are the level sets of f , with lighter colors corresponding to greater
values of f . The vectors represent the gradient vector field. The size of the arrows cor-
responds to the magnitude of the gradient at that point. Notice that the vectors are
perpendicular to the level sets.
The gradient plays a very important role in classical mechanics. For example, as we will
show in Chapter 5, for a particle with potential energy V (x, y, z), the force vector F (x, y, z)
acting on the particle is given by the negative gradient of the potential energy:
∂V ∂V ∂V
F (x, y, z) = −∇V (x, y, z) = − î − ĵ − k̂ (3.4.13)
∂x ∂y ∂z
Although we will discuss the relationship between force, energy, and equilibrium states of a
mechanical system in more detail in Chapter 5, it is worth noting that the above equation
says that forces point in the direction of most rapidly decreasing potential energy (because
of the minus sign).
i i
i i
i i
s1 = R[ 0 ] ∗R[ 1 ] ∗R[ 2 ]
print ( g r a d i e n t ( s1 , R) )
OUTPUT:
R y∗ R z ∗R. x + R x∗ R z ∗R. y + R x∗R y∗R. z
In the Mathematica code, we simply use the Grad[f, {x, y, z}] command. Notice that
the results are in a list, where the first element of the list is the x-component, the second
element is the y-component, and the third element is the z-component.
a 2î − 2ĵ + k̂ 2 2 1
â = q =p = î − ĵ + k̂
a2x + a2y + a2z 22 + 22 + (1)2 3 3 3
In order to get the form of the gradient in other coordinate systems, we return to the
fact that the total differential df of the scalar function is given by:
∂f ∂f ∂f ∂f ∂f ∂f
df = dx + dy + dz = dr + dθ + dz = ∇f · dr (3.4.14)
∂x ∂y ∂z ∂r ∂θ ∂z
where in Cartesian coordinates dr = dxî + dy ĵ + dz k̂. We can use this equation to obtain
expressions for the gradient in cylindrical and spherical coordinates, as follows.
The gradient of any scalar function f in cylindrical coordinates (r, θ, z) will be of the
general form:
∇f (x, y, z) = (∇f )r r̂ + (∇f )θ θ̂ + (∇f )z k̂ (3.4.15)
i i
i i
i i
where (∇f )r , (∇f )θ , (∇f )z are the components of the gradient along the unit vectors r̂, θ̂, ẑ
in cylindrical coordinates. From (3.3.19) we recall that dr in cylindrical coordinates is given
by:
dr = drr̂ + rdθθ̂ + dzẑ (3.4.16)
By evaluating the dot product of (3.4.15) and (3.4.16), and substituting into (3.4.14), we
obtain:
∂f ∂f ∂f
dr + dθ + dz = (∇f )r dr + (∇f )θ rdθ + (∇f )z dz (3.4.17)
∂r ∂θ ∂z
By equating the corresponding coefficients of (dr, dθ, dz) on the two sides of (3.4.17), we
find:
∂f 1 ∂f ∂f
(∇f )r = , (∇f )θ = , (∇f )z =
∂r r ∂θ ∂z
Therefore by using (3.4.15), the expression of the gradient in cylindrical coordinates is:
∂f 1 ∂f ∂f
∇f (r, θ, z) = r̂ + θ̂ + ẑ (3.4.18)
∂r r ∂θ ∂z
The gradient in spherical coordinates can also be obtained with the same method, and is
of the form:
∂f 1 ∂f 1 ∂f
∇f (r, θ, φ) = r̂ + θ̂ + φ̂ (3.4.19)
∂r r ∂θ r sin θ ∂φ
Note that our work above has also given us the del operator in various coordinate systems:
The “derivative” properties of the gradient are shown in the following identity, which is
similar to the derivative product rule:
i i
i i
i i
f = x∧ 2 + y ∧ 2 + z ∧ 2;
gradf = Grad[f, {x, y, z}]
gr1 = VectorPlot3D[gradf, {x, −2, 2}, {y, −2, 2}, {z, −2, 2}, VectorPoints → 5];
gr2 = Graphics3D[Sphere[{0, 0, 0}, 2]];
Show[gr1, gr2]
Out[ ]=
In vector calculus, the del operator is a vector differential operator. As we have seen,
when applied to a scalar field, the del operator produces the gradient of a scalar field ∇f .
However, when applied to a vector field, the del operator can produce two quantities, the
divergence of a vector field ∇ · v, and the curl of a vector field ∇ × v. In the next subsections,
we will illustrate these new quantities.
i i
i i
i i
If u, v are any vector functions and f is any scalar function, the divergence also obeys the
following product rule:
∇ · (f v) = f (∇ · v) + v · ∇f (3.4.26)
However, note that the divergence of a cross-vector product is less intuitive:
∇ · (u × v) = v · (∇ × u) − u · (∇ × v) (3.4.27)
OUTPUT:
R x∗R y + R x∗ R z + R y∗ R z
False
i i
i i
i i
Note that the divergence takes on a different form in cylindrical and spherical coordinates.
To compute the divergence in these coordinate systems, choose the appropriate form of the
del operator from (3.4.7), and compute the dot product of the appropriate form of the Del
Operator and the vector of interest. The results are (see also the problems at the end of
this chapter):
i i
i i
i i
1 1 1 ∂vr
∂ ∂vθ ∂
∇×v = vφ sin θ − r̂ +
− rvφ θ̂
r sin θ ∂θ ∂φ r sin θ ∂φ ∂r
1 ∂ (rvθ ) ∂vr
+ − φ̂ (3.4.37)
r ∂r ∂θ
The curl has many uses in physics. For example, as we will study in more detail in
Chapter 5, the necessary and sufficient condition for a force F (x, y, z) to be conservative is
that its curl equals zero:
The following are Python and Mathematica codes for the gradient.
We first compute the curl using the function curl from the Python SymPy library
package physics.vector. Notice that the curl is nonzero.
To check if a vector field is conservative we use the function is conservative which is
also from the SymPy library package physics.vector. In this example, the function returns
the value False, so that this field is not conservative.
v2 = R [ 1 ] ∗ R [ 1 ] ∗ R [ 2 ] ∗ R. x + R [ 0 ] ∗ R [ 2 ] ∗ R. y + R [ 0 ] ∗ R [ 1 ] ∗ R. z
print ( c u r l ( v2 , R) )
print ( i s c o n s e r v a t i v e ( v2 ) )
OUTPUT:
( R y ∗∗2 − R y ) ∗R. y + (−2∗R y∗ R z + R z ) ∗R. z
False
In Mathematica, we use the command Curl to find the curl of the vector function.
Notice that the curl is not equal to zero and, therefore, the vector field is not conservative.
i i
i i
i i
∂2 ∂2 ∂2
∇2 = ∇ · ∇ = + + (3.4.39)
∂x2 ∂y 2 ∂z 2
The Laplacian is ubiquitous throughout modern mathematical physics, appearing in
Laplace’s equation, Poisson’s equation, the heat equation, the wave equation, and the
Schrödinger equation—to name a few.
When the del operator acts on the gradient, divergence, and curl of a vector, it produces
five possible second derivatives; the use of the scalar Laplacian and vector Laplacian gives
two more. Here are these “second derivatives” created by applying the del operator.
In the above equations, we included the older notation of grad, div, curl, and ∆ with the
more modern del ∇ notation. You may encounter the older notation in some other books.
These second derivatives are of interest, principally because they are not always unique or
independent of each other. As long as the functions v and f are well-behaved, two of these
identities are always zero:
The first identity states that the curl of the gradient of any scalar field is always equal to
zero, and the second one states that the divergence of the curl of any vector field is always
equal to zero.
Two of the second derivatives are always equal:
∇ × (∇ × v) = ∇(∇ · v) − ∇2 v (3.4.50)
i i
i i
i i
∂2 ∂2 ∂2
2
+ 2 + 2 = ∇ · ∇ = ∇2
∂x ∂y ∂z
To find the curl of the gradient, calculate the cross product, (∇ × ∇V ). This quantity
is always zero, as seen in (3.4.47) for the curl of a gradient.
The dot product between two vectors A and B can be summarized by the following
relationships:
d dA dB
A · B = Ax Bx + Ay By + Az Bz = AB cos θ (A · B) = B+A
dt dt dt
The cross product between two vectors A and B can be summarized by the following
relationships:
a × b = |a||b| sin θ n̂
î ĵ k̂
a2 a3 a1 a3 a a2
a × b = a1 a3 = ĵ + 1
a2 î − k̂
b1 b2 b3 b1 b3 b1 b2
b2 b3
d da db
a × b = −b × a (a × b) = × b + a×
dt dt dt
There are several relationships that can simplify expressions with multiple dot or cross
products. They are:
i i
i i
i i
î ĵ k̂
∂ = ∂vz −
∂vy ∂vx ∂vz ∂vy ∂vx
∇ × v = ∂x + +
∂ ∂
∂y ∂z
î − ĵ − k̂
∂y ∂z ∂z ∂x ∂x ∂y
vx vy vz
The gradient gives the direction of most rapid change in a scalar field. The divergence
measures the magnitude of a vector field’s source or sink at a given point. The curl at a
point is proportional to the on-axis torque to which a tiny pinwheel is subjected.
The directional derivative df /dn in the direction n̂ = nx î + ny ĵ + nz k̂ is:
df ∂f ∂f ∂f
= n̂ · ∇f = nx + ny + nz
dn ∂x ∂y ∂z
A force F will be conservative if and only if its curl is zero: ∇×F = 0
Second derivatives of vector and scalar fields are also important in mechanics. One
such example, the Laplacian ∇2 is a scalar operator that can be applied to either vector
or scalar fields:
∂2 ∂2 ∂2
∇2 = ∆ = 2
+ 2+ 2
∂x ∂y ∂z
∇ × (∇f ) = 0
∇·∇×v = 0
Motion in two and three dimesions can sometimes be best described by non-Cartesian
coordinates. Those coordinates often involve using Euler’s identities:
x = r cos θ y = r sin θ
r = x2 + y 2 θ = tan−1 (y/x)
p
i i
i i
i i
i i
i i
i i
i i
i i
i i
13. By using the derivative properties of the dot and cross product, show that for any
vector v
dv d2 v dv d3 v
d
v. × 2 = v. × 3
dt dt dt dt dt
As an alternative method of proving this vector identity, use the symbolic capabilities
of Mathematica.
14. If the position of a particle on the xy-plane is given by r = cos(ωt)î + sin(ωt)ĵ , show
that the vector r × v is a constant vector. What is the physical meaning of this result?
15. Use the symbolic capabilities of Mathematica to prove the following identities for any
vectors A, B, C, and D:
(A × B) · (C × D) = (A · C) (B · D) − (A · D) (B · C)
A × (B × C) + B × (C × A) + C × (A × B) = 0
(A × B) · [(B × C) × (C × A)] = (A · [B × C])2
d2 r
m = f (r)r̂
dt2
where r is the position vector of P measured from an origin O, r̂ is a unit vector in
the direction r, and f (r) is a function of the distance of P from O.
(a) Show that r × dr
dt =c where c is a constant vector.
(b) Interpret physically the cases f (r) < 0 and f (r) > 0.
(c) Interpret the result in (a) geometrically.
(d) Describe how the results obtained relate to the motion of the planets in our solar
system.
dθ 2
h 2 i 2
a = dv
dt = d r
dt 2 − r dt r̂ + (r ddt2θ + 2 dr
dt dt )θ̂
dθ
i i
i i
i i
One possible method is by taking the time derivative of the velocity vector
dr dr dθ
v= = r̂ + r θ̂
dt dt dt
r = C θ̇−1
√
ṙ = A ln r + B
where A, B, C are constants.
(b) Find the velocity and acceleration of this point on the surface of the rolling wheel.
(c) The curvature k at a point of a curve is defined as
|a × v|
k=
|v|3
where a and v are the acceleration and velocity vectors of the wheel. Find the
curvature k at a point located at the top of the cycloid.
23. Find the radius of curvature at the top of the parabolic path of a projectile by assuming
no air resistance. The curvature k at a point of a curve is defined as
|a × v|
k=
|v|3
where a, v are the acceleration and velocity vectors of the object.
24. Find expressions for the velocity and the acceleration vectors in spherical coordinates.
i i
i i
i i
represent the coordinate system, we can define three vectors ei which are orthogonal
to each other by evaluating the partial derivatives:
∂r
ei = ei · ej = 0 (i = 1, 2, 3 i 6= j)
∂qi
These ei vectors can be normalized by dividing by their length, and the normalized
basis vectors are then:
ei ei
êi = =
hi |ei |
The lengths hi of the basis vectors êi are very useful functions known as scale factors
of the coordinates. In the case of Cartesian coordinates, qi = x, y, z and
∂r ∂r ∂r
êi = , , = î, ĵ, k̂
∂x ∂y ∂z
and it is easy to verify that hi = 1.The infinitesimal length in this orthogonal coordi-
nate system is found from:
√ q
ds = dr · dr = h21 dq12 + h22 dq22 + h23 dq32
The corresponding infinitesimal volume elements are calculated from:
dV = h1 h2 h3 dq1 dq2 dq3
(a) Find expressions for the unit vectors êi in cylindrical coordinates, and determine
the corresponding scale factors.
(b) Find expressions for the infinitesimal arc length ds2 and infinitesimal volume dV .
27. Repeat the previous problem for spherical coordinates, i.e., find the unit vectors êi ,
the square of the element of arc length ds2 , the corresponding scale factors and the
infinitesimal volume dV for spherical coordinates.
28. The gradient of a scalar field, f , the divergence and curl of a vector field
F = (F1 , F2 , F3 ), and the Laplacian of a scalar field f can be found using the
following general expressions:
i i
i i
i i
where the hi are the scale factors of the basis vectors, êi , which were defined in
Problems 26 and 27.
(a) Use the given expressions to obtain the gradient, divergence, and curl in cylindrical
coordinates.
(b) Obtain the divergence for spherical coordinates.
29. Consider the function V (x, y, z) = x4 + y 5 + z 6 in three dimensions. Calculate the
following quantities and indicate which of these is a vector and which is a scalar.
(a) ∇ ∇2 V
34. Prove: ∇ · (φA) = (∇φ) · A + φ (∇ · A), where φ is a scalar function and A is a vector
function.
36. If v = ω × r, prove ω = 1
2 curl v where ω is a constant vector.
i i
i i
i i
39. Maxwell’s famous equations concern the magnetic field vector B and the electric field
vector E. If these vectors satisfy the equations:
∂B ∂E
∇·E = 0 ∇ · B=0 ∇×E = − ∇×B =
dt dt
then show that E and B satisfy the generalized wave equation:
∂2E ∂2B
∇2 E = ∇2 B =
∂t2 ∂t2
40. Prove that for any vector field A, the following is true:
∇ × (∇ × A) = −∇2 A + ∇ (∇ · A)
i i
i i
i i
CHAPTER 4
Momentum, Angular
Momentum, and
Multiparticle Systems
In this chapter, we focus on two important physical quantities: linear momentum and angu-
lar momentum. Under certain conditions, as explained in this chapter, momentum and
angular momentum are conserved quantities, meaning that their values are constant during
a physical process. The conservation of momentum or angular momentum can then be used
to solve problems by comparing the state of a system before a physical process to the state
of that same system after this process. In this chapter, we will also study multiparticle sys-
tems, continuous mass distributions (i.e., systems that are not point particles), and center
of mass. We develop relationships between the total momentum of a system of particles to
the momentum of the system’s center of mass. We also develop a relationship between the
angular momentum of a system of particles to angular momentum of the system’s center of
mass. We will show that these relationships will also hold for continuous mass distributions.
Finally, we discuss numerical integration techniques which will be useful when finding a
system’s center of mass.
Newton’s third law will lead us to an interesting result. Consider two particles interacting via
a force, which we will refer to as an “internal force.” For example, if our two objects are the
Earth and Moon, their gravitational interaction would be the internal force. Furthermore,
assume that each particle is experiencing an external force due to other bodies outside of
the system. In the case of the Earth and the Moon, that could be the Sun’s gravitational
force, a body outside the Earth-Moon system. The net force F1 on particle 1 is then:
F1 = F12 + Fext
1 (4.1.2)
107
i i
i i
i i
where F12 is the force on Particle 1 exerted by Particle 2 (the “internal force”), and Fext 1
is the net external force acting on Particle 1. Similarly, Particle 2 experiences a net force:
F2 = F21 + Fext
2 (4.1.3)
Let p1 and p2 represent the linear momenta of the two particles, then Newton’s second law
states that
ṗ1 = F1 = F12 + Fext
1 (4.1.4)
and
ṗ2 = F2 = F21 + Fext
2 (4.1.5)
Now, we define the total momentum vector as P = p1 + p2 , to be the sum of the momenta
of each particle. Then,
Ṗ = ṗ1 + ṗ2 = Fext
1 + F2 =F
ext ext
(4.1.6)
Hence, we see that if there are no total net external forces Fext = Fext
1 + F2
ext then Ṗ = 0,
The total momentum of a system is conserved if no external forces act on that system. This
is an important law in physics, one that we will use to solve problems in cases where the
external forces add to zero.
Example 4.1: Inelastic collision between two bodies
Consider two objects with masses m1 and m2 and moving with velocities v1 and v2 ,
respectively. The two masses collide and stick together, moving away from the collision as
one object with a velocity v. Ignoring external forces occurring during the collision find
the velocity of the objects immediately after the collision.
Solution:
This type of collision is called a perfectly inelastic collision because the two bodies
are stuck together after the collision. Because we can ignore the external forces during
the collision, we can use conservation of momentum to solve this problem. The total
momentum before the collision is:
Pinitial = m1 v1 + m2 v2
Pfinal = (m1 + m2 ) v
Note that because the two objects stick together after the collision, both masses have the
same velocity, hence the right-hand side of Pfinal above. Conservation of momentum tells
us that Pinitial = Pfinal and therefore:
m1 v1 + m2 v2
v=
m1 + m2
Notice that the final velocity is the weighted average of the initial velocities, where the
masses of each object serve as the weights for the average. While this problem may seem
to be simple, it is an important application of the conservation of momentum; it is used
i i
i i
i i
v0 θ
m1 φ
m2
m2
v2
Solution:
You may remember from your introductory physics course that in an elastic collision
both the total momentum mv and total kinetic energy 1/2mv 2 are conserved. We will
discuss kinetic energy more in depth in the next chapter. For now, your knowledge from
your introductory physics course will suffice to solve this problem.
Only mass m1 is moving before the collision, and the problem states that the velocity
is in the x-direction. Therefore,
pinitial = m1 v0 î
The momentum after the collision involves both particles. By analyzing the velocity vectors
v1 and v2 into their x and y components. We can write:
m1 v0 = m1 v1 cos θ + m2 v2 cos φ
0 = m1 v1 sin θ − m2 v2 sin φ
1 1 1
m1 v02 = m1 v12 + m2 v22
2 2 2
The top two equations represent the total momentum in the x and y directions respec-
tively, and the final equation represents the conservation of kinetic energy. In addition,
we know that v0 = 10.0 m/s, m1 = 0.10 kg, and m2 = 0.20 kg. As you might imagine, this
problem can involve a lot of algebra. However, a CAS will be able to assist us with that
algebra. The following is an example of code written in Mathematica, which can be used
to solve this problem.
i i
i i
i i
θ = π/6;
v0 = 10;
m1 = 0.1;
m2 = 0.2;
OUTPUT: Solve::ifun: Inverse functions are being used by Solve, so some solutions
may not be found; use Reduce for complete solution information.
Next, we study the case where three or more particles are involved. The questions arises:
Do we still expect momentum to be conserved when there are many interacting particles?
Let’s begin by examining the net force acting on a collection of N interacting particles. The
net force on the ith particle is: X
Fi = Fext
i + Fij (4.1.7)
j6=i
i i
i i
i i
If this rewriting is not clear to you, set N = 3 and write out the terms in (4.1.8); you
will find that you can rewrite the sum as a collection of two terms, as shown previously.
Newton’s third law tells us that Fji = −Fij , and therefore the second term in the previous
equation is zero leading to:
XN
Ṗ = F = Fext
i (4.1.10)
i=1
In other words, the momentum of the system of particles is conserved if there are no external
forces acting on the system, like in the case of two-particle systems.
4.2 ROCKETS
A practical example of conservation of momentum is rocket propulsion. The rocket has a
challenge of moving forward without pushing against something. For example, an automobile
moves forward by turning its wheels so that it pushes against the ground. However, a rocket
in space has nothing to push against. Instead, the rocket solves this problem by ejecting
mass, similar to the recoil of a gun. Before a gun is fired, the bullet and the gun are at rest.
However, once the bullet is fired, it is propelled forward. In order for the total momentum to
be conserved, the gun has to move in the direction opposite that of the bullet. The velocities
of the gun and the bullet are, of course, very different because the gun has more mass than
the bullet does.
We begin by considering a rocket flying horizontally and assume that the rocket is not
experiencing any external forces. The rocket has momentum p(t) = mv at time t, where m
and v are the mass and velocity of the rocket. The momentum of the rocket is measured
by an observer in a reference frame that is at rest relative to the rocket. A short time later
denoted by t + dt, the rocket has expelled a small amount of mass dm0 , which is moving
at a speed u relative to the rocket, as illustrated in Figure 4.1. The speed u is sometimes
called the exhaust speed. The ejection of mass leads to an increase in the rocket’s speed by
a small amount dv, as measured in the rest frame. Therefore, the momentum of the system
immediately after the mass was ejected is:
where the second term above is the momentum of the ejected mass relative to the rest
frame.
dm' v+dv
m-dm'
u
Figure 4.1: A rocket after ejecting a small amount of mass dm0 at a speed relative to the
rocket, u.
i i
i i
i i
where the last line was obtained by noting that the change in mass of the rocket is
dm = −dm0 . Furthermore, we ignored the term dm0 dv because both dm0 and dv are small,
therefore their product is negligible. Dividing both sides by dt we obtain:
m v̇ = −u ṁ (4.2.6)
The left-hand side of (4.2.6) is Newton’s second law (F = ma), and therefore we see that
the net force acting on the rocket is equal to −u ṁ, which is sometimes called the thrust.
Note that ṁ < 0, so the thrust is positive (and points to the right in Figure 4.1). We can
also write (4.2.6) as:
dm
dv = −u (4.2.7)
m
This equation can be integrated by assuming a constant u. If at t = 0 the rocket’s speed is
v0 and its mass is m0 , we obtain:
m
0
v − v0 = u ln (4.2.8)
m
which gives the velocity of the rocket as a function of its mass m. Equation (4.2.8) tells us
that in order to make a rocket go as fast as possible, engineers need to create rockets with
large exhaust speed u, and with a large initial to final mass ratio m0 /m. The final speed of
the rocket will be obtained once all of the fuel is burned. Rockets with multiple stages that
are ejected can help maximize m0 /m, so that the rocket can achieve even greater speeds.
Next, we study the case of a rocket moving vertically, experiencing the force of gravity.
Our goal is to find an equation which describes the velocity of the rocket. We will use a
coordinate system where the positive y-direction is upward, in the opposite direction of the
force of gravity. In this case, gravity is an external force acting on the rocket, and therefore:
dp
= − mg (4.2.9)
dt
dp = − mgdt (4.2.10)
p(t + dt) − p(t) = − mgdt (4.2.11)
mdv + udm = − mgdt (4.2.12)
dv dm
m +u = − mg (4.2.13)
dt dt
In order to obtain an equation for v, we introduce the constant burn rate α:
dm
α=− (4.2.14)
dt
Using our definition of burn rate, the last line of (4.2.13) becomes:
i i
i i
i i
dv α
=−g+ u (4.2.15)
dt m
αu
dv = −g + dt (4.2.16)
g m
u
dv = − dm (4.2.17)
α m
Integrating the above equation using v(0) = v0 and m(0) = m0 produces:
g m
0
v − v0 = (m − m0 ) + u ln (4.2.18)
α m
This equation gives the velocity of the rocket as a function of mass m when the burn rate α
is constant. Sometimes (4.2.18) is rewritten by integrating the definition of the burn rate:
m − m0 = −αt. In that case, (4.2.18) becomes:
m
0
v − v0 = −gt + u ln (4.2.19)
m
The advantage of using (4.2.19) is that one can see two terms which affect the rocket’s
motion. The first term, −gt, is the standard kinematics term for a particle in free fall, and
represents the effect of gravity on the velocity. The second term, u ln(m0 /m), is the same
term we obtained for the horizontal motion of the rocket and represents the effect of thrust
on the rocket’s velocity.
N N N
1 X 1 X 1 X
X= mi xi Y = m i yi Z= mi zi (4.3.2)
M M M
i=1 i=1 i=1
Here ri P
is the position of the ith particle, which has mass mi . The total mass of the system
is M = mi . Equation (4.3.1) is actually three equations, one for each coordinate of R.
For example, in Cartesian coordinates the position of the ith particle is ri = xi î + yi ĵ + zi k̂
and we can write R = X î + Y ĵ + Z k̂.
The center of mass is the weighted average location of the mass in the system, with
the weights being the masses mi . It is a weighted average, similar to how grades are often
calculated in a course where maybe 50% of your final grade may be based on mid-semester
exams, 20% on homework assignments, and 30% on a final examination. The location of
the center of mass will be closer to the heavier particles. The next example illustrates this
point.
i i
i i
i i
If we want to extend (4.3.1) to continuous distributions of mass, i.e., objects that are not
well-described as point particles, then we need to make some changes. One way to approach
a continuous mass distribution is by thinking of it as a collection of infinitesimal masses
dm. Then the sum over discrete particles in (4.3.1) becomes an integral over infinitesimal
mass elements dm. The center of mass equations then become:
In the equations above, the total mass is found from M = dm, and like (4.3.1) equation
R
(4.3.3) actually consists of three equations, one for each coordinate. Typically, one does not
integrate (4.3.3) using the mass. Most often the integration is done over the length, area,
or volume of the object, when the object does not have a uniform density.
The following box summarizes how to calculate the total mass M by using the density
of the object in one-, two-, and three-dimensional situations. In the case of one-dimensional
objects, we use a linear mass density λ (mass per unit length in units of kg/m). In the
case of two-dimensional objects, we use a surface density σ (mass per unit area in units of
kg/m2 ), and for three-dimensional objects we use the familiar volume density ρ (mass per
unit volume in units of kg/m3 ).
i i
i i
i i
The next example will demonstrate how to apply (4.3.3) to a continuous mass distribu-
tion.
Example 4.4: Center of mass of an isosceles triangle
A lamina of uniform mass per unit area σ, is shaped into an isosceles
√ triangle shown
below. The triangle has two sides of length a and the base has length 2a. Find the center
of mass of this triangular lamina.
y
a a
dm
dy
dx
x
2a
Solution:
First note that due to the symmetry of the triangle, the center of mass lies along
the y-axis. Computing the center of mass of a mass distribution involves identifying two
components. The first component to identify is the mass element dm, shown in the diagram
of the triangle. Cartesian coordinates are a natural choice of coordinates for this problem.
Therefore, dm = σdxdy will be used for the mass element.
The second component needed to find the center of mass are the limits of the integral in
(4.3.3). The limits of the integral range over the size of the object. We included a horizontal
gray strip near the top of the triangle to illustrate that the width of the triangle varies
with the height, with the strip becoming narrower
√ as it moves up the triangle. The right
side of√the strip is bound by the line y = a/ 2 − x√and the left side is bound by the line
y = a/ 2 + x, because the triangle √ has a height a/ 2. Therefore, we can √ consider the left
end of the strip to be at x = y − a/ 2, and the right end to be at x = a/ 2 − y.
The
R y-coordinate of the center of mass can then be found by computing Y =
1/M ydm:
i i
i i
i i
√ √
1 y=a/ 2 Z a/ 2−y
Z
Y = √ yσdxdy
M y=0 x=y−a/ 2
Z a/√2 a/√2−y
σ
= yx dy
M 0 √
y−a/ 2
√
Z a/ 2
σ a
= 2y
√ − y dy
M 0 2
√
a/ 2
σ ay 2 2y 3
= √ −
M 2 3
0
a3 σ
= √
6 2M
To finish the problem, we need to find the total mass M .
√ √ √
y=a/ 2 Z a/ 2−y a/ 2
σa2
Z Z Z
a
M= dm = σdxdy = σ 2 √ − y dy =
y=0
√
x=y−a/ 2 0 2 2
a3 σ
OUTPUT: √
6 2m
h n o n oi
M = Integrate σ, y, 0, √a2 , x, y − √a2 , √a2 − y
a2 σ
OUTPUT: 2
Y /.m → M
OUTPUT: a
√
3 2
i i
i i
i i
z
a
x
Solution:
To solve this problem, we first need to identify a coordinate system. A cone can be
naturally described in cylindrical coordinates. Therefore we use dm = ρdV = ρrdrdθdz. By
symmetry, the center of mass is along the z−axis. The θ coordinate ranges from 0 to 2π,
and the z coordinate ranges from 0 to h. However, the radius of the cone depends on the
height of the mass element above the xy−plane. The equation of the edge of the cone in
the yz−plane in the figure is r = az/h, so that the radius r ranges from 0 to az/h. Putting
all of this together,
1
Z h Z 2π Z az/h
Z= zρrdrdθdz
M z=0 θ=0 r=0
az/h
2πρ h r2
Z
= z dz
M 0 2
0
2πρa2 h 3 πρa2 h2
Z
= 2 z dz =
2h M 0 4M
Solving for the total mass M :
Z h Z 2π Z az/h
M= ρrdrdθdz
0 0 0
az/h
h
r2
Z
= 2πρ dz
2
0
0
h
πρa2 πρa2 h
Z
= z 2 dz =
h2 0 3
i i
i i
i i
ρ∗z∗r
Z = Integrate , {θ, 0, 2π}, {z, 0, h}, {r, 0, az/h}
m
a2 h2 πρ
OUTPUT: 4m
OUTPUT: 31 a2 hπρ
Z/.m → M
OUTPUT: 3h
4
Notice that in order to perform the triple integral in cylindrical coordinates, we needed
to explicitly include the equation for r = az/h in the argument of the Integrate command.
f(x)
f(b)
f(a)
x
a b
Figure 4.2: The area under the curve f (x) between x = a and x = b can be approximated
using the trapezoidal rule, which estimates the area using the dashed trapezoid shown in
the figure.
The trapezoidal area under the curve between the points x = a and x = b is found by
taking the sum of the triangle and rectangle shown in Figure 4.2:
i i
i i
i i
Z b
f (x)dx = (b − a)f (a) + 1/2(b − a)(f (b) − f (a)) + O (b − a)3 f 00 (4.4.1)
a
b−a
=[f (b) − f (a)] + O (b − a)3 f 00 (4.4.2)
2
where the error term O (b − a)3 f 00 is not included in calculations. The f 00 in the term is
the second derivative of f (x) evaluated at a point between x = a and x = b that maximizes
the second derivative. This is done to give an upper bound on the error. As expected, the
trapezoidal rule is not a preferred algorithm when it comes to calculating integrals, because
the area of a more complicated function may not be well-approximated by a single trapezoid.
However, the trapezoidal rule is the basis of other algorithms, hence it is important to know.
by sampling the function f (x) at three points x0 , x1 , and x2 , equally spaced by a distance
h. The formula for Simpson’s rule is:
Z x2
h
f (x0 )dx0 = [f (x0 ) + 4f (x1 ) + f (x2 )] + O h5 f (4) (x0 ) (4.4.4)
x0 3
where h = x2 − x1 = x1 − x0 and x1 = (x0 + x2 )/2. The last term in (4.4.4) is the remainder
(or error) term and is not included in calculations. The fourth-order derivative is evaluated
at some value in the interval [x0 , x2 ]. While Simpson’s rule divides the area under f (x) into
more segments than the trapezoidal rule, even more segments are often required. In those
cases, the extended Simpson’s rule is needed.
The extended Simpson’s rule breaks the interval up into N equal size divisions using the
formula:
xN
1
Z
h
f (x0 )dx0 = (f1 + 4f2 + 2f3 + 4f4 + · · · + 2fN −2 + 4fN −1 + fN ) + O (4.4.5)
x1 3 N4
where fi = f (xi ), xi = x1 + ih, and h = (xN − x0 )/N . The error term at the end of the
extended Simpson’s rule is again not included in calculations. The material included in this
section demonstrates only two techniques for numerical integration. For more information
and more techniques, the interested reader should see [Press et al. (2007)].
In order to implement the extended Simpson’s rule, follow these steps:
1. Define the function to be integrated.
2. Define the lower and upper limits of the integral (a and b) and the stepsize h.
i i
i i
i i
3. Create an array whose elements are the terms in the sum of the right-hand side of
(4.4.5).
(a) Define the first element of the array to be f (a), and the last element to be f (b).
(b) Define the other elements as follows: the elements with an even index i are
evaluated using the expression 4f (a + ih), and elements with an odd index i are
evaluated using the expression 2f (a + ih).
4. Sum the array and multiply the result by h/3.
Algorithm 3 uses Python to demonstrate these three different algorithms: the trapezoidal
rule, the simple Simpson’s rule, and the extended Simpson’s rules, for this simple integral:
Z 3
x2 + 2x + 3 dx = 27 (4.4.6)
0
In Algorithm 3 we changed some notation for simplicity, and instead of using an array
for the extended Simpson’s rule, we used a for-loop to demonstrate an alternative method.
The variables a and b are the lower and upper limits of integration, respectively. Recall that
lines which begin with # are comments and are not executed by the Python interpreter.
Also recall that lines beginning with the word OUTPUT show the output of the code, and
are not part of the code. Because we are going to need to evaluate f (x) = x2 + 2x + 3 many
times, we begin the program by defining the function being integrated. Notice that the
trapezoidal and basic Simpson’s rules are simple enough to perform in one or two lines of
code each. Furthermore, we included the print command to show how formatting can be
done in Python.
The calculation of the extended Simpson’s rule is more involved than the other two
methods. The coefficients used in the extended Simpson’s rule depend on the index of each
term in the right-hand side of (4.4.5). The coefficient of fi is 4 when i is even, and is 2 when
i is odd, except for i = 1 and i = N . In order to handle these dependencies, we used the if
and elif conditional statements. The if statement tests whether or not a condition is true
and if so, the indented line is executed.
The first if statement tests to see if i = 1 or i = N . If so, then the term hf (x)/3 is
added to the sum. Notice that the variable extended contains the sum of the terms on the
right-hand side of (4.4.5), and that the command += adds to and overwrites the variable
extended. If i 6= 1 or N , then the next conditional statement is evaluated.
The next conditional statement is elif, or else-if, which tests to see if i is even by cal-
culating i modulo 2. In Python, the % symbol is used for the modulo calculation. The elif
statement provides an additional conditional statement to test after the first if statement
and is useful when there are multiple possible outcomes for a conditional (in this case i can
be even, odd, 1, or N ). Finally, if none of the above conditional statements return true,
then i must be odd. The else conditional statement tells Python what to execute if no other
prior conditionals are true; in this case the term 43 hf (x) is added to the sum.
Notice that the trapezoidal rule performs poorly in our example. However, the basic
Simpson’s rule (using two divisions) works better than the extended Simpson’s rule which
uses 10,000 divisions. This is a good example of higher-order calculations not always per-
forming better than lower-order ones. You should experiment with the number of divisions
used in Algorithm 3 to find how the accuracy of the calculation changes with the number
of divisions, N .
Next we will demonstrate how to use numerical integration algorithms to obtain the
center of mass of continuous mass distributions. In the following example, we will use
i i
i i
i i
def f (x ) :
r e t u r n x ∗∗2 + 2∗ x + 3
a = 0
b = 3
#t r a p e z o i d r u l e
t r a p e z o i d = ( f ( b)− f ( a ) ) ∗ ( b−a ) / 2 . 0
p r i n t ( ” The t r a p e z o i d r e s u l t i s : { : . 5 f } ” . format ( t r a p e z o i d ) )
OUTPUT: The t r a p e z o i d r e s u l t i s : 2 2 . 5 0 0 0 0
#simpon ’ s r u l e
c = ( a+b ) / 2
h = b − c
simpson = h ∗ ( f ( a ) / 3 . 0 + 4∗ f ( c )/3.0+ f ( b ) / 3 . 0 )
p r i n t ( ” The Simpson r e s u l t i s : { : . 5 f } ” . format ( simpson ) )
OUTPUT: The Simpson r e s u l t i s : 2 7 . 0 0 0 0 0
#extended simpson ’ s r u l e
N = 10000 #number o f s t e p s
extended = 0
h = ( b−a ) / (N−1) #need N s t e p s s t a r t i n g a t an i n d e x o f 0
f o r i i n r a n g e ( 0 ,N ) :
x = a + i ∗h
i f i == 0 o r i == N−1:
extended += h∗ f ( x ) / 3 . 0
e l i f i % 2 == 0 :
extended +=2.0∗h∗ f ( x ) / 3 . 0
else :
extended +=4.0∗h∗ f ( x ) / 3 . 0
i i
i i
i i
Y = Integrate 2 ∗ y 2 , y, 0, a 2 , x, y − a 2 ,a 2 −y mass
√ √ √
OUTPUT: 0.353553
X = Integrate 2 ∗ y ∗ x, y, 0, a 2 , x, y − a 2 ,a 2 −y mass
√ √ √
OUTPUT: 0.
Next we perform the integration using the numerical integration routine, nquad, from
SciPy’s integrate library. The command nquad uses the command quad over multiple
variables of integration. The command quad also comes from SciPy’s integrate library,
and uses a FORTRAN 77 library called QUADPACK [Piessens et al. (1983)] to evaluate
an integral. The code to perform the integration is shown in Algorithm 2.
Notice that we can include the bounds as a function; this allows us to use the non-
constant limits of integration necessary to describe the sides of the triangle. In addition
to finding Y , we also calculated X to show that X = 0, as expected.
i i
i i
i i
from s c i p y import i n t e g r a t e
import numpy a s np
a = 1.0
d e f sigma ( x , y ) :
return 2.0∗y
def fx (x , y ) :
r e t u r n x∗ sigma ( x , y )
def fy (x , y ) :
r e t u r n y∗ sigma ( x , y )
d e f y bounds ( ) :
r e t u r n [ 0 , 1 / np . s q r t ( 2 . 0 ) ]
d e f x bounds ( y ) :
r e t u r n [ y−1/np . s q r t ( 2 . 0 ) , 1 / np . s q r t (2.0) − y ]
OUTPUT
The x−c o o r d i n a t e i s : 0 . 0
The y−c o o r d i n a t e i s : 0 . 3 5 3 5 5 3 3 9 0 5 9 3 2 7 3 8 4
where we used the definition of the center of mass from (4.3.1). Therefore, the system of
particles moves like a single particle of mass M , acted upon by the external forces. We can
also compute the net momentum of the system of particles:
N N N
X X d X d
P= pi = mi ṙi = mi ri = M R = M Ṙ (4.5.2)
dt dt
i=1 i=1 i=1
i i
i i
i i
The total momentum of the system is equal to the momentum of the system’s center of
mass. Again, we can think of the system as a single particle of mass M located at the
system’s center of mass. The time derivative of the total momentum gives:
Ṗ = M R̈ = F (4.5.3)
In other words, the total linear momentum of the system is conserved if there are no net
external forces acting on the system.
In summary, we can think of the net external forces as acting on the center of mass of
the system. The results of this section also hold for continuous mass distributions, but the
summations would need to be replaced by integrals.
When it comes to translational motion, the above result allows us to think of the object
or system of particles as a point particle located at the object’s center of mass. The mass of
the point particle would equal the total mass of the object or the total mass of the system
of particles. Since the external forces act on the center of mass, we can simply follow the
motion of the object’s center of mass. For example, an American football is not a point
particle. However, if the football is punted, its center of mass would follow a parabolic
trajectory because the center of mass is being acted upon by the force of gravity, and obeys
the equations established for projectile motion in Chapter 3.
Solution:
The chain can be thought of as a system of particles. The force on the center of mass
of the chain is:
Fnet = M g − FN
where M g is the weight of the chain and FN is the force exerted by the table on the chain.
We have chosen a coordinate system where +x is in the downward direction, so that speed
of the falling chain is ẋ > 0. The mass of the falling part of the chain is ρ(a − x), because
the length of the chain that has yet to fall is (a − x). The momentum of the falling part
i i
i i
i i
Angular Momentum
` = r×p (4.6.1)
Because the position r of the particle is measured relative to the origin, the angular
momentum ` is also relative to the origin. The linear momentum p does not depend on the
origin, while by contrast the angular momentum vector ` depends on the choice of origin.
The angular momenta of multiple particles can only be compared (and added) if all angular
momenta are measured with respect to the same origin. Further, note that the particle does
not need to be revolving around the point O, in order for the particle to have an angular
momentum relative to O.
The cross product relationship between r and p tells us that the direction of angular
momentum is found using the right-hand rule. In this case, we can imagine taking the
fingers of our right hand and pointing them in the direction of r. Next, we can sweep (or
curl) our fingers in the direction N
of p resulting in the thumb, which points in the direction
of ` , pointing into the page. The symbol in Figure 4.3 represents the angular momentum
vector pointing into the page. If ` pointed out of the page, we would use a circle with a dot
inside to represent the direction of the vector.
If we compute the time derivative of the angular momentum, we obtain:
i i
i i
i i
r ℓ
O x
Figure 4.3: A particle of mass m at a location r relative to the origin O is moving with
momentum p, and has an angular momentum ` = r × p relative to the origin. In this case,
the angular momentum of the particle is into the page.
`˙ = r × F = N (4.6.3)
Therefore, the time rate of change of the angular momentum is the torque, N = r × F.
Like the angular momentum, the torque is measured relative to the origin, the same origin
as the angular momentum. Equation (4.6.3) is sometimes referred to as Newton’s second
law for rotation, even though rotational motion is not needed in order for the particle to
experience a torque, or to have angular momentum.
i=1
i i
i i
i i
ri'=ri-R
mi
R
ri
x
Figure 4.4: A collection of N particles. Each particle mi is at a location ri relative to the
origin and a position r0i relative to the center of mass. The empty circle is the center of mass
of the system.
N N
!
X X
mi r0i × Ṙ = mi r0i (4.7.4)
× Ṙ
i=1 i=1
N N
!
X d X
mi R × ṙ0i = R× mi r0i (4.7.5)
dt
i=1 i=1
i i
i i
i i
i=1
which says that the total angular momentum of the system is the sum of two terms:Pthe first
term is the sum of the angular momenta of the particles about the center of mass ( r0i × p0i
); the second term is the angular momentum of the center of mass about the origin (R × P).
If we compute the time derivative of L we find,
N
X
L̇ = ri × ṗi (4.7.8)
i=1
which is the net torque N acting on the system (recall that ṗ = F). Next, we will examine
the net torque acting on the system in detail:
N
X N
X
N= ri × ṗi = ri × Fi (4.7.9)
i=1 i=1
where ṗi = Fi is the net force acting on mass mi . This net force acting on mi is the sum of
external and internal forces, therefore:
N
X X
N= i +
ri × Fext Fij (4.7.10)
i=1 j6=i
N
X N X
X
= ri × Fext
i + (ri × Fij ) (4.7.11)
i=1 i=1 j6=i
i i
i i
i i
L̇ = Next (4.7.15)
Hence, the net angular momentum of the system is conserved only if the net external
torques are zero. The conservation of angular momentum can then be written as:
where we dropped the units in the determinant to simplify the equation. Next, we follow
the same steps for Particle 2:
î ĵ k̂
`2 = m2 (r2 × v2 ) = 1.0 3.0 2.0 0 = 6.0 kg m2 /s k̂
0.0 2.0 0
L = `1 + `2 = 2.0 kg m2 /s k̂
Note that we could add angular momenta, in this case, because each of the individual
angular momenta were calculated with respect to the same origin.
i i
i i
i i
The trapezoid rule and Simpson’s rule are used for numerically evaluating definite integrals.
For translational motion, we can think of a system of particles as a point particle
located at the object’s center of mass, and of the net external forces acting on the center
of mass of the system. If M is the total mass and P is the total momentum:
Ṗ = M R̈ = F
The angular momentum ` of a single particle is the cross product of position vector r and
linear momentum vector p:
` = r×p
and the torque N is the cross product of position vector r and force vector F:
`˙ = r × F = N
The angular momentum L of a system of particles is the sum of two terms, the angular
momenta of the particles about the center of mass, and that of the center of mass about
the origin (R × P):
XN
L= r0i × p0i + R × P
i=1
where R is the location of the center of mass relative to the origin, and P is the total
momentum. If we differentiate L with respect to time, we find Newton’s second law for
Rotation:
L̇ = Next
From Newton’s second law for rotation, we can write the law of conservation of angular
momentum for a system of particles:
If Next = 0, then L = constant
i i
i i
i i
i i
i i
i i
10. A model rocket has a burn rate α = −dm/dt. Using (4.2.19) derived in his chapter,
m
0
v − v0 = −gt + u ln
m
calculate the height the rocket as a function of time t.
11. A rocket has an initial mass m and burn rate α. What is the minimum exhaust velocity
u needed to allow the rocket to lift-off from the Earth’s surface immediately after its
engines fire?
12. Consider a two-stage rocket. During the first phase of flight, the rocket with mass m0
has exhaust velocity v. At the end of phase 1, the rocket has a mass m1 = mp + mf 1
where mp is the mass of the payload and mf 1 is the mass of the fuel container. At
burnout, the rocket ejects the mass mf 1 and begins to burn the fuel stored in the
payload. At the end of the second burnout, the rocket’s mass is m2 . The exhaust
velocity during the second burnout is also u. Find the velocity of the rocket after the
second burnout.
Sections 4.3–4.4: Center of Mass and Numerical Integration
13. Prove that the center of mass of a two-particle system lies along the line joining the
two particles.
14. Consider the cone in Example 4.5. Suppose it has a density of ρ = (3 + z)/2, where
the units of the constants are such that ρ is measured in kg/m3 . Find the center of
mass of the cone.
15. Compute the location of the center of mass of a uniform solid hemisphere of radius
R.
16. Find the center of mass of three particles located at r1 = 2i − 2j + 5k, r2 = 5i − 2k,
r3 = −i + j − k, with masses m1 = 2 m2 = 5 m3 .
17. Calculate the center of mass of a lamina with a density of σ = 2(x − 3)(y 2 − 7) formed
from the area defined by the curves y = x2 and y = x between x = 0 and x = 1.
18. A uniform piece of metal is cut into the shape of the function y = sin x from x = 0
to x = π. The metal has a density of σ = 2xy. Compute the location of the center of
mass.
19. Compute the location of the center of mass of a lamina cut in the shape of 1/3 of the
unit circle with a density of σ = 2xy.
20. Consider a system of three particles with masses and positions at time t: m1 = 1.0
is at r1 = 3t2 i + 2(t − 1)j − 3 cos (2t) k, m2 = 2.0 is at r2 = ti + 4t3 j − exp(−t)k, and
m3 = 1.5 is at r3 = i + t2 j − 3k. Each mass is measured in kilograms, each position is
measured in meters from the origin, and t is time measured in seconds. Find:
(a) The center of mass of the system.
(b) The total linear momentum of the system.
(c) The total angular momentum of the system.
(d) The net force acting on the system.
(e) The net torque acting on the system.
i i
i i
i i
28. Consider a rigid rotating object. By breaking the object up into many small pieces
of mass mi , show that the components of the object’s angular momentum parallel to
the axis of rotation can be written as L = Iω, such that:
X
I= mi ri2
where ri is the position of the mass element mi relative to the origin.
i i
i i
i i
29. Let L be the angular momentum of a system of particles relative to the origin O. Let
L0 be the angular momentum of the same system with respect to another origin O0 .
Find an equation for L0 in terms of L. You will need to consider the vector rO0 , the
position of O0 relative to O.
30. Prove that the total gravitational torque about the center of mass of a system of
particles is zero.
31. Consider a system of two particles experiencing the force of interaction:
f12 = k [(r2 − r1 )]
where k is a positive constant and f12 is the force Particle 1 exerts on Particle 2. What
is the net internal torque acting on the system? Discuss your result.
32. Consider a system of two particles experiencing the force of interaction:
i i
i i
i i
CHAPTER 5
Energy
In this chapter, we introduce the concepts of energy and conservative forces. We also discus
how energy can be used to describe the motion of a system, and as part of that discussion,
we introduce the concepts of equilibria stability. We begin by studying one-dimensional
systems because they are easy to visualize, and then extend the concept of work and energy
to two, and three-dimensional systems, including a discussion of line integrals. We conclude
the chapter with a discussion of the energy of multiparticle systems.
135
i i
i i
i i
where primes are used to distinguish between variables of integration and the limits of
integration. The unit of work, the Joule is 1 J = 1 N·m. Note that work is negative if the
object’s displacement is in the opposite direction of the force. For simplicity, we will focus
only on positive work right now.
Recall from Chapter 2 that if F = F (x), then we can write:
dv
F (x) = mv (5.1.2)
dx
which can be written as:
1
d
F (x) = mv 2 (5.1.3)
dx 2
Notice that the right-hand side of the above equation is the derivative of the quantity
T = 21 mv 2 , called the kinetic energy:
dT
F (x) = (5.1.4)
dx
Separating variables and integrating:
Z T Z x
dT 0 = F (x0 )dx0 (5.1.5)
T0 x0
Equation (5.1.6) is called the work-kinetic energy theorem, which states that work needs
to be done by the force F (x) in order to change the kinetic energy of a system. For example,
the force of gravity does work on a falling object, and its kinetic energy increases.
Additional insight into the nature of work can be obtained from (5.1.6). If the change
of the particle’s kinetic energy ∆T is positive, then positive work is being done by F and
the system is gaining kinetic energy. If however ∆T < 0, then F (x) is doing negative work
and the particle is losing kinetic energy.
Next, we will add another restriction to F (x). If the work done by the F (x) is indepen-
dent of the path taken by the particle from x0 to x, then we say that force is conservative.
In a later section, we will study forces in two and three dimensions and develop a deeper
mathematical formalism for conservative forces. However, in one dimension all forces of the
form F = F (x) are conservative. Let’s consider the force of gravity exerted by an object
near the Earth’s surface. In this case, the force is the weight W = −mg ĵ, where the negative
sign is introduced to denote that the weight points downward, in the negative y-direction.
The work done by gravity is,
Z y
W= −mgdy 0 = −mg(y − y0 ) (5.1.7)
y0
where we used y for vertical displacements. Notice that the work done by gravity depends
only on the end points of the particle’s motion.
i i
i i
i i
Energy 137
Consider a falling rock (in a vacuum, of course!). We noted earlier that the rock will
increase its kinetic energy as it falls. Where does that kinetic energy come from? According
to (5.1.6), the force of gravity does work on the rock and (5.1.7) tells us how much work
is done by gravity. Another way of thinking of the process is to say that during the fall,
the rock’s energy is changing from one form into another. While the kinetic energy is the
energy associated with the motion of the rock, there is an additional energy associated
with the configuration of the rock-Earth system. In this case, the important parameter
is the rock’s height above the ground. In fact, (5.1.7) supports the idea of configuration-
dependent energy. The term y − y0 in (5.1.7) says that the work done by gravity depends
on how the distance between the Earth’s surface and the rock changes during the motion.
That change of distance is a change in the configuration of the rock-Earth system. The
energy associated with the configuration of the rock-Earth system, or any system, is called
the potential energy. The potential energy function V (x) in one-dimension is defined by:
dV (x)
F (x) = − (5.1.8)
dx
Notice that (5.1.8) says that forces act so as to decrease the potential energy along x.
Hence, in our rock-Earth example, the force of gravity is decreasing the rock’s potential
energy and increasing the rock’s kinetic energy. Furthermore, this example illustrates that
forces transfer one form of energy into another, in this case, potential energy transferring
to kinetic.
To find the potential energy for a given force, F(x), we need to integrate the force:
Z x
V (x) − V (x0 ) = − F (x0 )dx0 (5.1.9)
x0
This equation is known as the conservation of mechanical energy, stating that the sum
of the kinetic and potential energies, called the total mechanical energy, is always constant
if only conservative forces act on the system. Conservative forces conserve energy!
Conservative forces are very useful in physics as they allow us to work with a scalar
quantity, the mechanical energy E = T + V . Mechanical energy is a useful way of solving
physics problems because instead of needing to keep track of the magnitude and direction
i i
i i
i i
of vectors, as one needs to do with forces, all one needs to do is keep track of how the value
of the different types of energies are changing as the particle moves. Looking at (5.1.12), we
see that solving problems using conservation of energy requires only that we compare one
state of the system to another. For example, suppose we drop a rock starting at rest from
a height of 5 m above the ground and we want to know the kinetic energy of the rock after
it fell 2 m. All we need is to know is the initial kinetic energy (which is zero) and the initial
and final potential energies in order to obtain the final kinetic energy (and thus the speed).
We don’t need information about the state of the system during the actual fall!
Equation (5.1.12) also tells us that if the potential energy of the system decreases, its
kinetic energy must increase and vice versa, in order to maintain the equality of both sides
of the equation. The ability to interpret the meanings behind equations such as (5.1.8) and
(5.1.12) is an important skill for physicists.
T0 + V (x0 ) =T + V (x)
1 1 1
0 + kA2 = mv 2 + kx2
2 2 2
Therefore, using:
1 1 1
E = mv 2 + kx2 = kA2 (5.1.14)
2 2 2
we find that the total mechanical energy in terms of A is E = 21 kA2 .
(b) At the turning points xTURN the velocity vTURN = 0 so (5.1.14) gives:
1 1
E = 0 + kxTURN = kA2
2 2
Therefore, the expression for the turning points in the motion of a simple harmonic
oscillator is:
xTURN = ±A = 2E/k (5.1.15)
p
i i
i i
i i
Energy 139
where we used Wnc to denote that the work is done by a nonconservative force. Note that
we explicitly included the minus sign in (5.1.16) because the force of friction opposes the
direction of motion, and we consider the displacement to be in the positive direction.
Consider the case of a particle, again constrained to move along a line, experiencing
multiple forces, both conservative and nonconservative. The net work Wnet is the total
work done by all of the forces. The net work can be written as Wnet = Wc + Wnc , where Wc
is the work done by the conservative forces and Wnc is the work done by nonconservative
forces. In regards to the work above, Wnet = ∆T and Wc = −∆V . Therefore,
∆T =Wnc − ∆V (5.1.17)
Wnc =∆T + ∆V (5.1.18)
Wnc =∆E (5.1.19)
or, in other words, the work done by nonconservative forces is equal to the change of the
system’s mechanical energy. So in cases where conservative forces act on the system, kinetic
energy gets exchanged for potential energy and vice versa. However, when nonconservative
forces act on the system, then mechanical energy can be lost or gained. An object falling
through air loses potential energy to both kinetic energy and work done by friction heating
the object and the surrounding air. However, if we account for the heat, we will find that
the total energy is conserved.
In the next section, we will study how to use potential energy to describe the motion of
an object.
Hence, the local maxima or minima of the potential energy gives the location of the system’s
equilibrium points, x0 . As we will see, the plot of V (x) can be used to describe the motion
of a particle.
Physical systems can be in three types of equilibria: stable, unstable, and neutral. An
example of a stable equilibrium is shown at the position near x = 0.7 mm in Figure 5.1.
This corresponds to values for which the potential energy V (x) is a local minimum. In
i i
i i
i i
other words, d2 V /dx2 > 0 at the equilibrium position. In the case of stable equilibrium,
the forces acting on the particle will tend to restore the equilibrium. The restoring nature
of the force can be seen in the graph of V (x), by looking at the slope of V (x) on either
side of the equilibrium point. Note that the sign of the force is opposite of the sign of the
potential energy’s slope. Hence a displacement to the right of the equilibrium, there V (x)
has a positive slope, is met with a negative (left pointing) force. Similarly, to the left of
the equilibrium, V (x) has a negative slope, and therefore a particle displaced to the left of
equilibrium experiences a positive (right pointing) force.
An unstable equilibrium is shown at x = 0 in Figure 5.1. This corresponds to values for
which the potential energy V (x) is a local maximum, or in other words, d2 V /dx2 < 0 at
the equilibrium position. This means that if the particle is displaced an arbitrarily small
distance from the equilibrium state, the force causes it to move even farther away. By
studying the slope of V (x) on either side of the equilibrium, we can see that the particle
experiences a positive (rightward) force when displaced to the right of equilibrium, and a
negative (leftward) force when displaced to the left of the equilibrium position, hence the
force acts to move the particle farther from equilibrium.
1.00
0.95
Potential V(x)
0.90
0.85
0.80
0.75
-1.0 -0.5 0.0 0.5 1.0
position x (mm)
Figure 5.1: A graph of a potential energy function which shows stable equilibrium points
near x = 0.7 and x = −0.7 and an unstable equilibrium at the origin (x = 0).
A final type of equilibrium not shown in Figure 5.1 is called neutral equilibrium. A neutral
equilibrium corresponds to situations for which the potential energy V (x) is a constant, and
hence its derivative is zero. In the case of neutral equilibrium, the system will tend to remain
in equilibrium if displaced by a small amount.
i i
i i
i i
Energy 141
CM
h-R
h
H
R θ R
Solution:
The figure on the left shows the object in its equilibrium state (upright). The figure
on the right shows the object tilted an angle θ from the vertical. The vertical line passing
through the center (circle) and center of mass (square) denotes the equilibrium position.
From the second diagram we can see that the height of the center of mass is H = R + (h −
R) cos θ. Therefore, the potential energy is:
V = mg [R + (h − R) cos θ] (5.2.2)
Next, to get the stability of the equilibrium state, we take two derivatives of V :
dV
= − mg(h − R) sin θ (5.2.3)
dθ
d2 V
= − mg(h − R) cos θ (5.2.4)
dθ2
Notice that the first derivative dV /dθ correctly identifies θ = 0 as an equilibrium con-
dition, since sin θ = 0 at θ = 0. Furthermore, if θ = 0 is inserted into the second derivative,
we find that the second derivative is negative when h > R. Therefore, the equilibrium at
θ = 0 is unstable when h > R. Additionally, the upright position is stable if h < R, since
in that case the second derivative is positive.
i i
i i
i i
70
60
Potential V(x)
50
40 x3 E2
Out[ ]=
30
x1 x2
20
E1
10
0.0 0.5 1.0 1.5 2.0 2.5 3.0
Position x
Figure 5.2: A potential energy function demonstrating turning points.
energy. Past the equilibrium point, the force begins to oppose the particle’s velocity (note
that V (x) has a positive slope in this region). Because the potential energy is increasing, the
particle’s kinetic energy would decrease in order to conserve the total mechanical energy. At
x = x2 , the particle has lost all of its kinetic energy and comes to a stop. However, it is still
experiencing a leftward force, and so the particle begins to move back towards x1 , where it
will once again come to a stop and move back towards x2 . We see that the resulting motion
is an oscillation between the turning points, x1 and x2 . Both methods of description, the
“frictionless track” and the more detailed analysis in this paragraph, tell us that a particle
with an initial position of x0 = x1 and initial velocity v0 = 0, will oscillate about the stable
equilibrium reaching an amplitude of x1 to the left of equilibrium and x2 to the right of
equilibrium.
We can, of course, use mathematics to describe this motion. Using energy conservation,
E = T + V = mv 2 /2 + V (x), we can solve for the velocity as a function of position:
2
r
v=± (E − V (x)) (5.2.5)
m
For a particle with energy E = E1 , we can see from Figure 5.2 that V (x1 ) = E1 and V (x2 ) =
E1 and therefore, using (5.2.5), the velocity is zero at the turning points, x = x1 and x = x2 .
Furthermore, for x > x2 and x < x1 , the velocity is a complex number because V (x) > E in
those regions. Complex velocities are not physical, and this means that the particle does not
have enough energy to access those regions. For a particle with E = E1 , the regions x < x1
and x > x2 are called classically forbidden regions. According to the physics of classical
mechanics, the particle cannot be in classically forbidden regions because it doesn’t have
enough energy to access those regions. However, such regions are accessible in quantum
mechanical systems.
Using v = dx/dt, we can find the position by integrating (5.2.5):
Z x
dx0
t − t0 = q (5.2.6)
x0 2
m (E − V (x 0 ))
Equation (5.2.6) can be a difficult way to solve for x(t) in closed form, for all but the most
simple of potential energy functions.
i i
i i
i i
Energy 143
We can repeat our qualitative analysis for the motion of a particle starting at rest at the
unstable equilibrium x3 in Figure 5.2. In that case, if the particle is perturbed to the right,
it increases its velocity as the potential decreases. The graph ends at x = 3, so we cannot
tell what will happen to the particle for x > 3. If the particle is deflected to the left, then it
will increase in speed until it passes through the stable equilibrium, after which it will slow
down until it reaches x ≈ 0.1, before turning around. Mathematically, the particle will then
reach the unstable equilibrium coming to a stop.
Notice that both Figures 5.1 and 5.2 have local minima. The motion around local min-
ima corresponds to an oscillatory motion, if the energy of the particle is low enough. Let
us understand why that is the case using mathematics. Consider a particle near a stable
equilibrium x = x0 , hence, the equilibrium position x0 needs to be a local minimum of
V (x). Next, we ask the question, what is the mathematical form of V (x) near the stable
equilibrium point? In order to obtain the mathematical form of V (x) near a stable equi-
librium, x0 , we will need to perform a Taylor expansion in terms of x around an energy
minimum (x = x0 ). A Taylor expansion gives a polynomial approximation of a function near
a particular point. The general form of the Taylor expansion is:
1
V (x) = V (x0 ) + (x − x0 )V 0 (x0 ) + (x − x0 )2 V (2) (x0 ) + O(x − x0 )3 (5.2.7)
2
where O(x − x0 )3 indicates terms of the order of (x − x0 )3 and higher, and the notation
V 0 (x0 ) and V (2) (x0 ) indicate the first and second derivatives of V evaluated at x = x0 ,
respectively. If we consider the case where the displacement from equilibrium is small, then
x − x0 is small and terms O(x − x0 )3 and higher are negligible. Because V (x0 ) is a minimum,
V 0 (x0 ) = 0, so the linear term drops out:
1
V (x) = V (x0 ) + (x − x0 )2 V (2) (x0 ) + O(x − x0 )3 (5.2.8)
2
The constant term V (x0 ) is arbitrary and thus may also be dropped. This is equivalent to
saying that we can set the zero point of the energy to be any value, hence we set V (x0 ) = 0.
Or to put it another way, constant energy terms do not affect the particle’s motion, because
the force acting on the particle is the derivative of the potential energy, and therefore
constant terms do not contribute to the force. With all of this taken into account, we obtain
the form of V (x) for the simple harmonic oscillator :
1 1
V (x) ≈ (x − x0 )2 V (2) (x0 ) = k(x − x0 )2 (5.2.9)
2 2
where we set k = V (2) (x0 ). Hence, we find that Hooke’s Law appears any time we are
studying small amplitude oscillations about a stable equilibrium. Thus, given an arbitrary
potential energy function V (x) with a local minimum (and a non-vanishing positive second
derivative), we can use the solution to the simple harmonic oscillator found in Example
2.7 to provide a solution for small perturbations around the equilibrium point. As the
perturbations get larger, however, we need to include higher order terms in the Taylor
expansion of V (x), and Hooke’s Law no longer applies. These so-called nonlinear oscillators
can be very difficult to work with. We will study nonlinear oscillators in Chapter 13.
It is worth mentioning that (5.2.9) gives us a method of finding the frequency of small
oscillations about the stable equilibrium. Recall thatpthe frequency of oscillations for a
particle trapped in a Hooke’s Law potential is ω = k/m. Therefore, the frequency of
small oscillations about a stable equilibrium x0 can be found from:
r r
k V (2) (x0 )
ω= = (5.2.10)
m m
i i
i i
i i
where σ, , x, and V are in appropriate SI units. Setting σ = 1 and = 1/4, determine the
type of equilibrium possible in this potential, and obtain a parabolic approximation of
this potential near the equilibrium point.
Solution:
We find the equilibrium point by setting the derivative of V (x) equal to zero and
solving for x:
dV (x)
dx = 12x−13 − 6x−7 = 0
x = 21/6 = 1.122
We can easily verify that this is a stable equilibrium point by examining the second
derivative. We can use Mathematica to expand the function V (x) around the equilibrium
point, and plot both the series approximation and the potential. The command NSolve is
used to find the equilibrium point x0 = 1.122, and the command Series is used to obtain
the first two terms of the series expansion:
V (x) = −0.25 + 7.143(x − x0 )2
V = 1
x12
− x16 ;
i i
i i
i i
Energy 145
0.4
Lenard-Jones
Approximation
potential, V(x)
0.2
Out[ ]=
0.0
-0.2
In the plot above, the dashed line is the series approximation near the stable equilib-
rium, and the solid line is the Lenard-Jones potential. Notice that the approximation fails
to provide a reliable estimate of the potential energy the farther the particle moves from
the equilibrium position.
The code above introduces some new Mathematica commands. The double ampersand,
&& is Mathematica’s AND operator. It tells Mathematica to solve for when V 0 (x) = 0
AND x > 0. The commands Normal and Chop, are used to turn the series expansion into
a result that can be used for further calculations.
dV a
F (x) = − = 2 − 2bx + 3cx2
dx x
The next step is to use Newton’s second law in order to create a second-order ODE,
which can be solved for x(t):
a
mẍ = 2 − 2bx + 3cx2
x
This can be reduced to two first-order equations using the variable, y ≡ ẋ:
ẋ =y
1 a
ẏ = − 2bx + 3cx2
m x2
Next, we insert the known values of a, b, and c, and m in order to use Euler’s method
to solve for x(t) and y(t) for the given initial condition: x(0) = 0.21 and y(0) = 0. Below
is a Python implementation, where the algorithm needed to generate the graph of x(t) is
removed for brevity.
i i
i i
i i
m, a , b , c = 1 , 4 , 2 5 , 8
x0 = 0 . 2 1 #i n i t i a l p o s i t i o n
y0 = 0 #i n i t i a l v e l o c i t y
def force (x ) :
r e t u r n a/x ∗∗2 − 2∗b∗x + 3∗ c ∗x ∗∗2
x , y = [ x0 ] , [ y0 ]
dt = 0 . 0 0 0 1
t = np . a r a n g e ( 0 , 2 , dt )
f o r i i n r a n g e ( 1 , l e n ( t ) ) : #Euler ’ s method
x . append ( x [ i −1] + dt ∗y [ i −1])
y . append ( y [ i −1] + dt ∗ f o r c e ( x [ i − 1 ] ) )
1.0
0.9
0.8
0.7
position (m)
0.6
0.5
0.4
0.3
0.2
0.0 0.5 1.0 1.5 2.0
time (s)
The graph of x(t) is below the Python program. Notice that the motion is periodic,
oscillating between points x = 0.21 and x = 0.95, as suggested in Figure 5.2.
Notice the dot product between the force F and the infinitesimal displacement dr in
(5.3.1). The dot product tells us that work calculations involve only the component of the
force parallel to the displacement, i.e., the component of the force perpendicular to the
displacement does no work. Note that (5.3.1) is also applicable to one-dimensional systems,
and we no longer need the restriction that F is in the same direction as dr.
i i
i i
i i
Energy 147
The right hand side of (5.3.1) is a line integral. In other words, the integral is evaluated
along a specific path. The value of W can depend on the path taken by the particle. It is
common to rewrite (5.3.1) including the path C followed by the particle:
Z
W= F · dr (5.3.2)
C
In order to calculate integrals of this type, one can break them up into component integrals,
by writing out the dot product of the components of the force vector F = Fx î + Fy ĵ + Fz k̂
and the differential dr =dxî + dy ĵ + dz k̂, as follows:
Z Z Z
W= Fx dx + Fy dy + Fy dy (5.3.3)
C C C
and then evaluate the individual integrals along the path C. The next several examples
demonstrate how to calculate path integrals.
Example 5.5: Calculating the work done by a force along a straight line path
Find the work done by the force F = xy î − y 2 ĵ along the straight line path connecting
the origin to the point (2, 1) as indicated in the figure below.
(2,1)
y
y=x/2
(0,0) x
Solution:
This problem is most easily done in Cartesian coordinates. We write for the work
Z Z Z Z Z
W= Fx dx + Fy dy + Fz dz = xydx − y 2 dy (5.3.4)
C C C C C
The straight line path shown in the figure can be written as y = x/2, and therefore
the differentials are related by dy = dx/2. The x-coordinate along path AB changes from
x = 0 to x = 2. By substituting these into the above expression, we end up with simple
integrals over the x-axis only:
x=2 x=2 2
3 2
Z Z Z
x x 2 dx
W= x dx − = x dx = 1 (5.3.5)
x=0 2 x=0 2 2 0 8
Example 5.6: Calculating the work done by a force along a path consisting of
multiple straight-line segments
Find the work done by the force in the previous example along the line path (0, 0) →
(2, 0) → (2, 1), indicated in the figure that follows. Does the work done by the force depend
on the path taken between points (0, 0) and (2, 1)?
i i
i i
i i
(2,1)
y
(0,0) (2,0)
x
Solution:
The work in this case is calculated in two segments, work as the sum of two segments
(0, 0) → (2, 0), added to the work in segment (2, 0) → (2, 1).
In the first segment we have y = 0 and therefore dy = 0, while x varies from x = 0 to
x = 2.
Z Z
W= Fx dx + Fy dy (5.3.6)
(0,0)→(2,0) (0,0)→(2,0)
Z Z
= xydx − y 2 dy (5.3.7)
(0,0)→(2,0) (0,0)→(2,0)
= 0+0 (5.3.8)
=0 (5.3.9)
The first integral above is zero because y = 0 along the path from (0, 0) → (2, 0). The
second integral is zero because dy = 0 along the path.
Along the second segment from (2, 0) → (2, 1), we have x = 2 and dx = 0, while y varies
from y = 0 to y = 1.
Z Z
W= Fx dx + Fy dy (5.3.10)
(2,0)→(2,1) (2,0)→(2,1)
Z y=1 Z y=1
= xydx − y 2 dy (5.3.11)
y=0 y=0
1
1 3
= 0− y (5.3.12)
3 0
1
=− (5.3.13)
3
The first integral above is zero because dx = 0 along this segment.
Examples 5.5 and 5.6 show clearly that the work done by this force depends on the
path taken between the two points (0, 0) and (2, 1). Such a force is called a nonconservative
force; we will see later on that nonconservative forces have no potential energy associated
with them.
Example 5.7: Calculating the work done by a force along a curved path
Find the work done by the force in Example 5.5 along the unit quarter circle from the
point (1, 0) to the point (0, 1) as indicated in the following figure.
i i
i i
i i
Energy 149
y
(0,1)
θ
x
(1,0)
Solution:
This is a problem that is best done in polar coordinates. Let θ be the angle measured
counterclockwise with respect to the x-axis, denoted in the figure above. A particle moving
along the quarter circle would have a constant radial coordinate r = 1. Therefore, the
conversion from Cartesian to polar coordinates would be: x = cos θ and y = sin θ. In order
to rewrite the force vector, we need to convert the unit vectors from Cartesian to polar
coordinates. To derive the conversion, we start with the unit vectors for polar coordinates
from Chapter 3:
r̂ = cos θî + sin θĵ
θ̂ = − sin θî + cos θĵ
Next, we multiply r̂ by cos θ and θ̂ by sin θ and subtract. The result is:
Next, we multiply the r̂-equation by sin θ and the θ̂-equation by cos θ, and add. The
result is:
ĵ = sin θr̂ + cos θθ̂
By inserting the coordinate transformations for x, y, î, and ĵ into the force we obtain:
i i
i i
i i
dT = F · (vdt) (5.4.3)
dT = F · dr (5.4.4)
where T0 is the value of the kinetic energy at position r0 . Equation (5.4.5) is the work-kinetic
energy theorem and is similar to (5.1.6) derived for one-dimensional systems.
where we have set a reference point r0 such that V (r0 ) = 0. The above definition works for
our purposes because we have a conservative force. If the force was not conservative, then
the integral in (5.5.1) would also depend on the path taken by the particle and would not
be simply a function of r.
Now suppose that a conservative force, F(r), moves a particle from a point r0 to the
point r2 along two different paths. One path goes directly from the reference point r0 to
r2 , while the other path goes from r0 to r1 and then from r1 to r2 . Because the work is
conservative, the work done along the two paths is the same:
Or, the work done by the conservative force F(r) is related to the potential energy
through the formula:
Combining equation (5.5.5) with (5.5.1), we see that the work done by a conservative force
along a closed-loop path must be equal to zero, since the change of potential energy along
that path is zero.
i i
i i
i i
Energy 151
and therefore:
∆(T + V ) = 0 (5.5.9)
which tells us that the total mechanical energy E = T + V remains constant as the particle’s
position changes from r1 to r2 .
In the case where nonconservative forces work on the system, (5.1.19) still applies for
two- and three-dimensional systems. However, in this case, (5.4.1) is used to compute T ,
(5.5.1) is used to compute V , and the work done by the nonconservative force Fnc (r) is
found using: Z r2
Wnc = Fnc (r0 ) · dr0 (5.5.10)
r1
We have seen that the potential energy V (x) can be found from the conservative force F,
by using (5.5.1). Now suppose we know V and want to find F. We know from our discussion
of one-dimensional systems that F = −dV /dx; however, how does that generalize to higher
dimensions? To derive the relationship between F and V in three dimensions, we begin by
considering an infinitesimal displacement r → r + dr over which a conservative force F(r) is
acting on a particle. The work done by that force is:
dW = F(r) · dr = Fx dx + Fy dy + Fz dz (5.5.11)
where we are working in Cartesian coordinates for simplicity, but any other coordinate
system can be used. Continuing with the same displacement r → r + dr, the work done by
a conservative force can also be computed as:
dW = − dV (5.5.12)
= − [V (r + dr) − V (r)] (5.5.13)
∂V ∂V ∂V
= − V (r) + dx + dy + dz + · · · − V (r) (5.5.14)
∂x ∂y ∂z
∂V ∂V ∂V
=− dx + dy + dz (5.5.15)
∂x ∂y ∂z
where r = xî + y ĵ + z k̂ , and in (5.5.15) we kept only linear terms in the multivariable Taylor
series. Comparing the coefficients of dx, dy, and dz in Equations (5.5.11) and (5.5.15), we
find:
∂V ∂V ∂V
Fx = − Fy = − Fz = − (5.5.16)
∂x ∂y ∂z
In other words, we can equate the first equation in (5.5.11) with the last equation in (5.5.15)
to show that the force is related to the potential energy through the gradient.
i i
i i
i i
Figure 5.3: Surfaces of constant gravitational potential energy and the direction of gravita-
tional force, Fg .
Notice the minus sign in (5.5.17). The gradient points in the direction of the most rapid
increase in a scalar function. In the case of (5.5.17), we see that forces point in the direction
of the most rapidly decreasing potential energy. Consider the gravitational potential energy
as illustrated in Figure 5.3. Points at the same altitude all have the same potential energy.
We can think of the dashed lines in Figure 5.3 as surfaces of constant potential energy.
As the altitude increases, the potential energy increases. Hence, the direction of the most
rapid increase in V is upward i.e., perpendicular to the ground (shown as the rectangle).
Because Fg = −∇V, the force vector points along the direction of the most rapid decrease
in V , which is downward.
A conservative force F can be written as the gradient of the potential energy. From
Chapter 3, we recall that the curl of a gradient of any scalar function V is always zero:
∇ × F = −∇ × ∇V = 0 (5.5.19)
Therefore, one method to identify whether or not a force is conservative is to compute its
curl. A force F will be conservative if and only if its curl is zero:
Conservative Force
∇×F = 0 (5.5.20)
i i
i i
i i
Energy 153
Solution:
Calculate the curl of the force:
∂Fz ∂Fy ∂Fx ∂Fz ∂Fy ∂Fx
∇×F = − î + − ĵ + − k̂
∂y ∂z ∂z ∂x ∂x ∂y
∇ × F = (0 − 0) î + (0 − 0) ĵ + (0 − x) k̂ = − xk̂ 6= 0
Example 5.9: Proving a force is conservative and finding the associated poten-
tial energy
Consider the force F = −k y î + xĵ − z k̂ .
(a) Show that it is conservative.
(b) Find the potential energy associated with this force.
Solution:
î ĵ k̂ h i
∇ × F = ∂x (−k) = + 0 + (1 1) =0
∂ ∂ ∂
∂y ∂z −k î(0) ĵ k̂ −
y x −z
So F is conservative and there exists a potential energy V associated with the force F.
Z Z Z Z
V = − F · dr = − Fx dx − Fy dy − Fy dy
Z Z Z
V =k ydx + k xdy − zdz
We must have
∂V ∂V ∂V
F = −∇V = − ,− ,−
∂x ∂y ∂z
So
∂V
Fx = −
= −ky
∂x
Integrating with respect to x while keeping y, z as constants:
i i
i i
i i
ri'=ri-R
mi
R
ri
x
Figure 5.4: A collection of discrete particles. The center of mass is located at R and is
denoted by the empty circle.
i i
i i
i i
Energy 155
First, we will calculate the kinetic energy of the system, by using a method similar to
our calculation of the total momentum of a system of N particles in Chapter 4. The total
kinetic energy is:
N
1
X
T= 2
mi vi (5.6.1)
2
i=1
By using ri = R + r0i
and ṙi = Ṙ + ṙ0i
(recall that a dot above a variable means differentiation
of that variable with respect to time), we can write vi2 as:
2
vi2 = ṙi · ṙi = ṙ0i + Ṙ · ṙ0i + Ṙ = vi0 + 2 ṙ0i · Ṙ + V2 (5.6.2)
where vi0 = ṙ0i is the velocity of the ith particle relative to the center of mass, and V = Ṙ is
the velocity of the center of mass. Inserting (5.6.2) into (5.6.1) gives:
N N
! N
1 X 1
X
0 2 d X 0
T= + Ṙ · mi ri + 2
(5.6.3)
mi vi mi V
2 dt 2
i=1 i=1 i=1
The middle term in the above equation contains mi r0i , which in Chapter 4 was shown to
P
be zero. Hence, the total kinetic energy becomes:
The first term in (5.6.4) is the kinetic energy of the center of mass, and the second term
in (5.6.4) corresponds to the kinetic energy of each particle that is moving relative to the
center of mass.
For potential energy, we need to consider both external and internal forces and their
corresponding potential energies. The potential energy of the particle with mass mi is:
X
Vi = Viext + int
Vi,j (5.6.5)
j6=i
where Viext is the potential energy due to the external force acting on mi , and Vi,j int is
the potential energy due to the interaction force between mi and mj . Examples of such
interaction forces could be gravitational,
electrostatic, or any other central forces. In general
for central forces we have Vi,j
int = V int r − r , meaning that the potential energy depends
i,j i j
on the relative distance between the two particles. Therefore the total potential energy of
the system is:
XN N X
X
V = Viext + int
(5.6.6)
Vi,j ri − rj
i=1 i=1 j6=i
In the case of rigid bodies the particles are at fixed distances, and therefore the internal
potential energy can be ignored, since it will be constant for each particle. Hence, for a rigid
body we will only have to worry about external forces. We will study rigid bodies in more
detail in Chapter 11.
i i
i i
i i
The Work-Kinetic Energy Theorem states that work needs to be done by the force
F (x) in order to change the kinetic energy of a system.
Zx
1 1
mv 2 − mvo2 = F (x0 )dx0
2 2
xo
The potential energy function V (x) associated with the force F (x) in one dimension is
defined by:
Z x
V (x) − V (x0 ) = − F (x0 )dx0
x0
The force F (x) associated with a potential energy V (x) can be found through differenti-
ation:
dV (x)
F (x) = −
dx
A particle is in an equilibrium state at x = x0 if:
dV
=0
dx
x0
Physical systems can exist in three types of equilibria: stable, unstable, and neutral
depending on the sign of d2 V /dx2 at the equilibrium point x = x0 .
The velocity is zero at the turning points of motion, so all of the particle’s energy is in
the form of potential energy.
Hooke’s law appears any time we are studying small amplitude oscillations about a
stable equilibrium, by doing a Taylor expansion of the potential V (x):
1 1
V (x) ≈ (x − x0 )2 V (2) (x0 ) = k(x − x0 )2
2 2
In two and three dimensions, the work done by the force F moving an object between
points r1 and r2 is: Z r2
W (r1 → r2 ) = F · dr
r1
i i
i i
i i
Energy 157
In order to calculate integrals of this type, one can break them up into component integrals:
Z Z Z
W= Fx dx + Fy dy + Fz dz
C C C
E = T + V (x) = To + V (xo )
The kinetic energy of a system of N particles is the sum of the kinetic energy of the
center of mass as well as the kinetic energy of each particle that is moving relative to the
center of mass:
N
1 1
0 2
X
T = MV + 2
mi vi
2 2
i=1
In a rigid body the particles are at fixed distances, and the internal potential energy can
be ignored, so that we only have to worry about external forces.
i i
i i
i i
5. A particle with mass m experiences a force F , and has velocity v = a/x (a > 0) along
the +x-direction. How much work is done by the force moving the particle from
position x1 to x2 , where x2 > x1 ?
i i
i i
i i
Energy 159
−c2 (x2 + k 2 )
V (x) =
x4 + 3c4
where c > 0 and k > 0. Find the equilibrium points as a function of x/c. Are the points
stable or unstable? Plot V (x) and discuss the possible types of motion of a particle
for different values of k/c.
19. A particle of mass m is acted upon by a one-dimensional potential energy V (x) =
3/4kx4 . The particle oscillates between two turning points, x1 and x2 . Find the period
of oscillation of a 1 kg particle in the case where E = 10 J and k = 3.0 N/m4 . Plot the
period of oscillation as a function of total energy and discuss your results.
Section 5.3: Work and Line Integrals
20. Compute the work done by the force F = xy î − y ĵ along the path joining the origin to
the point (2,4) along the following paths:
i i
i i
i i
(a) Along the x-axis to the the point (2,0), then parallel to the y-axis to the point
(2,4).
(b) Along the straight line path connecting the origin to the point (2,4).
(c) Along the path y = x2 .
21. Compute the work for the force F = y î − xĵ along the following closed-loop paths:
(a) The unit square in the first quadrant of the Cartesian plane, with one corner at
the origin.
(b) The unit square whose center is at the origin.
(c) The unit circle whose center is at the origin.
2
x2
(d) The ellipse 4 + y25 = 1.
Section 5.5: Conservative Forces and Potential Energy
22. Which of the forces below are conservative? For each conservative force, find the
potential energy.
(a) F = y î − xĵ
(b) F = −kr
(c) F = r sin (θ) eiφ r̂
(d) F = kxyzr̂
(e) F = k
r3
r̂
23. Prove that any central isotropic force F = F (r)r̂, where r2 = x2 + y 2 + z 2 , is conser-
vative.
24. We showed that a conservative force can be written as, F(r) = −∇V (r), and that
when conservative forces act on a system, the total mechanical energy is conserved
(∆(T + V ) = 0). Now consider a time-dependent force written as F(r, t) = −∇V (r, t).
Is energy conserved for this system? If so, prove it. If not, how is the conservation of
energy equation ∆(T + V ) = 0, changed?
(b) V = −kr cos θ, where k is a positive constant and r and θ are polar coordinates.
(c) V = −kr cos φ + cr2 sin θ, where k and c are positive constants and r, θ, and φ
are the typical spherical coordinates.
i i
i i
i i
Energy 161
29. Consider three particles. Particle 1 has a mass of m1 = 3.0 kg and is located at
r1 = 3.0t2 î at time t. Particle 2 has a mass of m2 = 1.0 kg and is located at r2 = 3.0t2 k̂.
Particle 3 has a mass of m3 = 2.0 kg and is located at r3 = 4.0tî − 8.0t2 ĵ + 3.0t3 k̂. In
each case, the position of the particle is measured in meters. Find:
(a) The location of the center of mass for the system.
(b) The kinetic energy of the center of mass.
(c) The kinetic energy of the system of particles by finding the individual kinetic
energies and adding them.
30. Consider a planet of mass M which is orbited by a Moon of mass m1 = m. The Moon
is in a circular orbit with an orbital radius of a. A dwarf planet, of mass m2 = m and
with kinetic energy of T2 approaches the two objects. The dwarf planet gets trapped
in the same circular orbit as the Moon, and the Moon is kicked free. Write down the
total energy of the three particle system (there will be two kinetic energies and three
potential energies). What is the kinetic energy of the Moon after the collision, after
it is far away? Note that the gravitational potential energy can be found in Problem
6 of this chapter.
31. Consider a group of three particles: m1 = 3.0 kg at r1 = 3î − 2ĵ + 3k̂, m2 = 2.0 kg at
r2 = −î + 3k̂, and m3 = 1.0 kg at r3 = 7ĵ − 5k̂, where each position is measured in
meters. If the particles are located far away from any other mass, compute the total
gravitational potential energy for this particle configuration.
32. The virial theorem states that for a gravitationally bound distribution of masses, the
average kinetic energy of the system is equal to minus one-half of the total potential
energy, hT i = − 12 hV i, where h i denotes an average. Prove that the virial theorem
holds for a single particle of mass m in a circular orbit around another particle of
mass M .
i i
i i
i i
CHAPTER 6
Harmonic Oscillations
As discussed in Chapter 5, when a system is displaced only slightly from a stable equilibrium
position, it oscillates with harmonic motion about this equilibrium position. In this chapter,
we study the problem of harmonic motion. We first introduce general concepts of linear
ordinary differential equations and present the simplest application of these equations for a
linear harmonic oscillator, with and without damping forces. This is followed by the study
of forced or driven oscillations under an external force, and the important concepts of
amplitude resonance, energy resonance, and the associated Q-factor for oscillatory systems.
The chapter will conclude by introducing the physical concepts of phase space and the
superposition principle, together with the important mathematical technique of Fourier
analysis.
dn x dn−1 x dx
An + An−1 + .. + A1 + A0 x = f (t) (6.1.1)
dtn dtn−1 dt
where x = x(t) measures the oscillator’s displacement from a stable equilibrium, A0 , . . . , An
are constants, and f (t) is a function of time t. The order of a differential equation is the
highest derivative in the equation, and the equation is called homogeneous when the function
f (t) = 0. Equation (6.1.1) is called a linear differential equation because it consists only
of terms that are linear in x(t) and their derivatives. In general, equations that describe
oscillations are nonlinear in x and its derivatives. However, as we will see later in this
chapter, in the limit of small displacements from equilibrium, the equations of motion
usually take a linear form similar to that of (6.1.1).
In this chapter, we will first study homogeneous linear second-order differential equations
of the form:
163
i i
i i
i i
d2 x dx
A2 + A1 + A0 x = 0 (6.1.2)
dt2 dt
If x1 (t), x2 (t) are two solutions of this equation, then the general solution xc (t) of (6.1.2)
is a linear combination of these two solutions, in the form:
d2 x dx
A2 2
+ A1 + A0 x = f (t) (6.1.4)
dt dt
In the general theory of differential equations, the solution of this equation is given by
the sum of the solution of the homogeneous equation xc (t) from (6.1.3), sometimes called
the complementary solution, and a particular solution xp (t) of the full nonhomogeneous
equation. The particular solution xp (t) is any solution of the equation:
d2 xp dxp
A2 2
+ A1 + A0 xp = f (t) (6.1.5)
dt dt
The general solution x(t) of a nonhomogeneous second-order differential equation then takes
the form:
General Solution of Nonhomogeneous Equation
i i
i i
i i
For small displacements, we can expand F (x) as a Taylor series about the equilibrium
position x = 0:
dF
F (x) ≈F (0) + x + O(x2 ) (6.2.1)
dx
x=0
We know that F (0) = 0, because the restoring force is not exerted when the system is at
equilibrium. Likewise, we know that the first derivative of F (x) with respect to x must be
negative, in order for F (x) to point towards equilibrium. In other words, the sign of F (x)
is opposite that of x. Therefore, we define the spring constant as:
dF
k≡− (6.2.2)
dx
x=0
We can ignore second-order and higher terms in (6.2.1) because x is small. Therefore, the
result is that for small displacements, the restoring force takes on the form of Hooke’s law:
Hooke’s Law
If x(t) is the location of the mass at a time t, and we assume that the restoring force
takes the form of Hooke’s law (i.e., small displacements from equilibrium), then we can
write Newton’s second law, F = ma = −kx(t), as:
d2 x
+ ω02 x = 0 (6.2.4)
dt2
i i
i i
i i
where ω02 ≡ k/m is a constant. Equation (6.2.4) is the differential equation of motion for the
so-called simple harmonic oscillator (SHO). The SHO is an important equation in physics
because it is the generic equation of motion for an object experiencing small amplitude
oscillations about a stable equilibrium. Simple harmonic motion results when a particle
experiences only a restoring force that is proportional to the particle’s displacement from
equilibrium, i.e., when there are no drag forces or other external forces acting on the particle.
To find an analytical solution to (6.2.4), we need a function x(t) whose second derivative
gives itself back after differentiation, with a multiplicative constant and a minus sign. In
other words, x(t) must satisfy: ẍ = −ω02 x. Recall that time derivatives can be represented
by dots above a function. While one obvious choice is either a sine or a cosine function, we
will find that an exponential solution will be more useful (and still equivalent). We look
for solutions that have the form x = Ceλt , where C is a constant. By substituting this into
(6.2.4) we find:
λ2 Ceλt + ω02 Ceλt = 0 (6.2.5)
or by simplifying:
λ2 + ω02 = 0 (6.2.6)
or q
λ± = −ω02 = ±i ω0 (6.2.7)
and thus λ must be one of the complex numbers i ω0 or −i ω0 . The general solution, then,
will be a linear combination of the two solutions:
x(t) = C1 [cos ( ω0 t) + i sin (ω0 t)] + C2 [cos (ω0 t) − i sin (ω0 t)] (6.2.9)
p
with C = A2 + B 2 tan φ = B/A (6.2.12)
i i
i i
i i
Returning to Equations (6.2.10) and (6.2.11), A, B, and C are now real constants which
can be determined by the initial conditions, which specify the state of the system at a given
time (usually at t = 0). For example, given the initial conditions v(0) = v0 and x(0) = x0 at
t0 = 0, we then have from (6.2.10) with t = 0:
and so A = x0 .
The speed of the mass v(t) is found by taking the derivative of (6.2.10) with respect to
time:
dx
v(t) = = −Aω0 sin ( ω0 t) + B ω0 cos ( ω0 t) (6.2.15)
dt
By using the given initial condition v(0) = v0 at t0 = 0 :
and so B = v0 /ω0 . Therefore the complete solution for these initial conditions is:
v0
x(t) = x0 cos ( ω0 t) + sin ( ω0 t) (6.2.17)
ω0
The above solution is the same result we obtained in Chapter 2, by substituting a
solution of the form x = A cos(ωt − φ) in the original equation mẍ = −kx.
i i
i i
i i
We can calculate the average kinetic energy < T > in one period of oscillation by taking
a time average of T over the period τ = 2π/ω0 as follows:
1 τ1 1 τ1 1
Z Z
<T >= mv 2 dt = mA2 ω02 cos2 (ω0 t)dt = mA2 ω02 (6.2.20)
τ 0 2 τ 0 2 4
Rτ
where we used the integral value τ1 0 cos2 (ω0 t)dt = 21 .
Similarly, we can find the average potential energy < V > in one period of oscillation by
taking a time average of V (t) :
1 τ1 2 1 τ1 2 2 1 1
Z Z
<V >= kx dt = kA sin (ω0 t)dt = kA2 = mA2 ω02 (6.2.21)
τ 0 2 τ 0 2 4 4
In conclusion:
1 1 1
<V >=<T >= mA2 ω02 = kA2 = E (6.2.22)
4 4 2
where we used k = mω 2 and E = 21 kA2 . This equation shows that the time averages of both
the kinetic and potential energies within one period of oscillation are equal to one-half the
total energy E.
This equation is an example of the well-known virial theorem in mechanics. Since the
potential and kinetic energies are also functions of the position x, we can also evaluate their
space averages over a period. The result for these space averages is:
i i
i i
i i
We can find the equation of motion for the simple plane pendulum by applying Newton’s
second law to the tangential axis only, F = −mg sin θ = ma and a = −g sin θ. The linear
2
acceleration a = ddt2s , where s is the arc length, is related to the angle θ by the formula
d2 s
s = Lθ and is given by dt2
= Lθ̈. Therefore we obtain:
By substituting ω02 = g/L, we obtain the equation of motion for the simple plane pendulum:
For small displacements from equilibrium, we can approximate sin θ ≈ θ, and the equation
becomes:
Equation of Motion for Small Pendulum Oscillations
θ̈ + ω02 θ = 0 (6.2.27)
where φ is a constant phase angle value, dependent on the initial conditions. The period of
oscillation for small displacements is found using τ = 2π/ω0 and is:
i i
i i
i i
Notice that the period of oscillation is independent of the amplitude of oscillation. This
condition holds only for small amplitude oscillations. Once the small angle approximation
sin θ ≈ θ no longer holds, then the period of oscillation will be related to the amplitude of
oscillation. Such a relationship is common in nonlinear oscillators, which we will study in
Chapter 13.
The gravitational potential energy of the pendulum is given by the expression:
V = mgh = mgL(1 − cos θ) (6.2.30)
Notice that the above expression is chosen such that V = 0 when θ = 0. In other words, the
potential energy is zero at the stable equilibrium point. By using the small angle approxi-
mation cos θ ' 1 − θ2 /2, we obtain:
1
V = mgLθ2 (6.2.31)
2
The total energy of the pendulum is then:
2 2
1 ds 1 1 dθ 1
E= m + mgLθ2 = mL2 + mgLθ2 (6.2.32)
2 dt 2 2 dt 2
where s = Lθ was used. Next, we can use (6.2.28) to find the total energy of the pendulum:
1
E = mgLθ02 (6.2.33)
2
This equation is of course completely analogous to the equation E = 12 kA2 for a mass-spring
oscillator. Notice that mgL appears in (6.2.33) instead of k. The restoring torque acting
on the pendulum is N = −mgL sin θ ≈ −mgLθ and is linearly proportional to the angular
displacement θ of the pendulum (for small displacements), with a coefficient of k = mgL.
Therefore (6.2.33) is consistent with the results for the spring-mass oscillating system.
i i
i i
i i
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
#d e f i n e i n i t i a l c o n d i t i o n s and a r r a y s
x = np . z e r o s ( num steps )
v = np . z e r o s ( num steps )
T = np . z e r o s ( num steps )
V = np . z e r o s ( num steps )
x[0] = 2
v[0] = 0
T[ 0 ] = 1/2∗m∗v [ 0 ] ∗ v [ 0 ]
V[ 0 ] = 1 / 2 ∗ k∗x [ 0 ] ∗ x [ 0 ]
f i g , a x e s = p l t . s u b p l o t s ( 1 , 2 , f i g s i z e =(10 ,4) )
a x e s [ 0 ] . p l o t ( time , x )
a x e s [ 0 ] . s e t x l a b e l ( ” time , s ” )
a x e s [ 0 ] . s e t y l a b e l ( ”x , m” )
a x e s [ 1 ] . p l o t ( time , T, ” r−−” , time , V, ”b” )
a x e s [ 1 ] . s e t x l a b e l ( ” time , s ” )
a x e s [ 1 ] . s e t y l a b e l ( ”V−e n e r g y=Blue , T−e n e r g y=Red , J ” )
f i g . s a v e f i g ( ” ex6 −1. e p s ” )
i i
i i
i i
2.0 2.00
1.5 1.75
V-energy=Blue,T-energy=Red, J
1.0 1.50
0.5 1.25
x, m
0.0 1.00
−0.5 0.75
−1.0 0.50
−1.5 0.25
−2.0 0.00
0 2 4 6 8 10 0 2 4 6 8 10
time, s time, s
Figure 6.3: The Python code output, using the Euler method to integrate the equation for
simple harmonic motion.
ẍ + 2γ ẋ + ω02 x = 0 (6.4.3)
where we used again ω0 = k/m for the natural frequency of the undamped oscillator.
p
Once more we look for solutions that have the form Ceλt , where C is a constant. By
substituting this into (6.4.3) and dividing all terms by Ceλt , we find:
i i
i i
i i
This implies a pair of solutions eλt , corresponding to the two roots of the quadratic. The
general solution is a linear combination of these two solutions:
p p
−γ+ γ 2 −ω02 t −γ− γ 2 −ω02 t
p p
−γt γ 2 −ω02 t − γ 2 −ω02 t
x(t) = A0 e + A1 e =e A0 e + A1 e
q (6.4.6)
or by introducing the parameter ω ≡ γ 2 − ω02 :
q q
ω= γ 2 − ω02 = (b/2m)2 − (k/m) (6.4.8)
The exact mathematical form of this equation x(t) q depends on the numerical values of
γ and ω0 , which determine whether the parameter ω = γ 2 − ω02 is real, imaginary, or zero.
Therefore there are three physically different behaviors of the damped harmonic oscillator as
follows: overdamped oscillations (γ > ω0 ), underdamped oscillations (γ < ω0 ), and critically
damped oscillations (γ = ω0 ).
two real exponential functions. The resulting motion is that the system returns to equilib-
rium by an exponential decay without oscillating. Alternatively, one can write the solution
as a linear combination of hyperbolic functions (see Problem 24, in the end of chapter
problems for this section):
The constants A0 , A1 , B0 , B1 can be found from the initial conditions of the oscillator.
Example 6.2 illustrates the motion of an overdamped oscillator using Mathematica to solve
for the equations of motion in a mass-spring system exhibiting overdamped behavior.
i i
i i
i i
m = 1; k = 1; b = 3; xo = 1;
Plot[x1, {t, 0, 15}, FrameLabel → {“Time t,s”, “Position x(t)”}, Frame → True,
LabelStyle → (FontSize → 20), ImageSize → Medium]
1.0
Position x(t)
0.8
0.6
0.4
0.2
0.0
0 2 4 6 8 10 12 14
Time t,s
The solution of the underdamped case oscillates repeatedly through the equilibrium
point as the amplitude of oscillation decays to zero. This is unlike the overdamped case,
where the oscillator may not pass through equilibrium at all, as its amplitude decays to
zero. This is an important difference in the behaviors of the two oscillators. Also notice from
i i
i i
i i
the underdamped oscillator, which is smaller than the natural frequency ωo = k/m.
p
The expression Ae−γt in (6.4.12) and (6.4.13) represents the decrease of the amplitude
as a function of time t for underdamped motion. Notice that the amplitude of oscillation
decays exponentially with time. The quantity τA = 1/γ has dimensions of time, in order for
the argument of the exponential e−γt to be dimensionless, and characterizes how fast the
amplitude of oscillation reaches zero. As a general rule of thumb for exponential decaying
functions, the amplitude of oscillation is considered to have reached a value of effectively
zero after five characteristic times, i.e., after t ∼
= 5/γ.
For small damping, i.e., when when γ << ω0 , we can use the binomial approximation
∼ 1 − ax to obtain:
(1 − x)a =
γ2
q q
2∼
ωd = ω0 − γ = ω0 1 − (γ/ω0 ) = ω0 1 − 2
2 2 γ << ω0 (6.4.15)
2ω0
Example 6.3 gives a Mathematica program for a mass-spring system following under-
damped behavior.
m = 1; k = 1; b = 0.1; xo = 1;
Plot[x1, {t, 0, 70}, FrameLabel → {“Time t,s”, “Position x(t)”}, Frame → True,
LabelStyle → (FontSize → 20), ImageSize → Medium]
i i
i i
i i
1.0
Position x(t)
0.5
0.0
-0.5
0 10 20 30 40 50 60 70
Time t,s
where A and B are determined by the initial conditions of the system. A comparison of x(t)
for the underdamped, overdamped, and critically damped oscillations is shown in Figure
6.4, where we can easily see that the critically damped oscillator (dash dot curve) returns
to equilibrium in a shorter time than other cases. Note that the same parameter values are
used for m, k, x(0), and v(0) to produce all three curves, while the value of b is variable.
1.0
Underdamped
Overdamped
0.5 Critical
Position x(t)
0.0
Out[ ]=
-0.5
0 2 4 6 8 10 12 14
Time t,s
Figure 6.4: Dependence of the harmonic oscillator behavior on the amount of damping, show-
ing underdamped (oscillatory curve), overdamped cases (solid line), and critically damped
(dashed curve). The code which produces these graphs is given in Examples 6.2, 6.3, and
6.4.
Example 6.4 gives a Mathematica program for a mass-spring system following critically
damped behavior.
i i
i i
i i
Plot[x1, {t, 0, 15}, FrameLabel → {“Time t,s”, “Position x(t)”}, Frame → True,
LabelStyle → (FontSize → 20), ImageSize → Medium, PlotStyle → {Thickness[.01]},
PlotRange → All]
1.0
Position x(t)
0.8
0.6
0.4
0.2
0.0
0 2 4 6 8 10 12 14
Time t,s
i i
i i
i i
these approximations the term γ sin (ωd t) is much smaller that the term ωd cos (ωd t) and
can be dropped in (6.5.3), to obtain:
dx
= A0 e−γt [ωd cos (ωd t) ] (6.5.4)
dt
Substituting into (6.5.1) we obtain:
1 2 1 2
E = m A0 e−γt ωd cos (ωd t) + k A0 sin (ωd t) e−γt (6.5.5)
2 2
By collecting terms:
1
E = A20 e−2γt mωd2 cos2 (ωd t) + k sin2 (ωd t) (6.5.6)
2
Substituting ωd ' ω0 = k/m we obtain:
p
1 1
E = A20 e−2γt k cos2 (ωd t) + k sin2 (ωd t) = kA20 e−2γt (6.5.7)
2 2
The above equation shows that the energy decays exponentially with time. The char-
acteristic decay time is τE = 1/(2γ). We previously saw in (6.4.12) that the corresponding
characteristic decay time for the amplitude is τA = 1/γ.
In summary, the total energy E and the rate of decrease of the total energy dE/dt for
the cases of weak damping γ << ωd , are given by:
The total energy lost within a small time interval ∆t can be estimated from:
dE
∆E ∼
= ∆t = −2γE∆t (6.5.10)
dt
where we used (6.5.9). It is customary to define the dimensionless Quality factor or Q-factor
Q as the ratio Q = 2π|E/∆E|, for a time interval ∆t equal to the period τ = 2π/ωd . By
substituting this value of ∆t = τ = 2π/ωd in this definition of Q:
E 1 ωd
Q = 2π = 2π = (6.5.11)
∆E 2γτ 2γ
In section 6.7, we will provide a more in-depth discussion of the quality factor, in con-
nection with the concept of energy resonance of a forced harmonic oscillator.
i i
i i
i i
where γ = b/2m, ω02 = k/m, and D = F0 /m. You might be wondering why we limit ourselves
to studying driving forces in the form of cosine functions. Clearly not all driving forces in
nature are in the form of cosine or sine functions. As we will show later in this chapter,
any periodic force (under certain conditions, which are often met in physical problems) can
be approximated using a series of sines and cosines. Therefore, the case of cosine or sine
driving forces is applicable widely and for many physical systems.
Note that (6.6.1) is a nonhomogeneous ordinary differential equation. As we discussed in
the introduction of this chapter on the general theory of differential equations, the solution
of nonhomogeneous equations is given by the sum of two parts, the homogeneous solution
xc (t), plus the particular solution xp (t) of the full nonhomogeneous equation:
We already saw in the previous section that the general solution xc (t) of the homogeneous
equations is given by (6.4.7):
q
xc (t) = e−γt A0 eωt + A1 e−ωt ω = (b/2m)2 − (k/m) (6.6.3)
The solution xc (t) is commonly referred to as the transient solution, because the exponential
decay term e−γt causes xc (t) to decay to zero. After the transient solution has decayed to
zero, the only solution left is xp (t) (assuming that xp 6= 0), and is therefore known as
the steady-state solution. Many physical systems exhibit both a transient and steady-state
behavior. Physicists are often (but not always!) more interested in the steady-state behavior
of the system, because that is the system’s long-term behavior.
Before we study the solution x(t) for the full equation (6.6.1), it is instructive to study
the case in which damping is absent. When no friction is present, the driven oscillator is
described by b = 0, so we obtain:
We try a solution of the form x = A cos (ωt − φ), and after we substitute and collect terms
we obtain:
−Aω 2 cos (ωt − φ) + ω02 A cos (ωt − φ) =D cos(ωt)
(6.6.5)
A ω02 − ω 2 cos (ωt − φ) =D cos(ωt)
This will have two solutions, when φ = 0 and cos (ωt − φ) = cos (ωt), and when φ = π where
cos (ωt − φ) = − cos (ωt). The amplitude A must be always positive, so we can write the two
i i
i i
i i
D
A= φ = π and ω > ω0 (6.6.7)
ω 2 − ω02
D
x = A cos (ωt) = cos (ωt) φ = π and ω > ω0 (6.6.9)
ω 2 − ω02
A plot of the amplitude A and of the phase difference φ as a function of the external
frequency ω is shown in Figure 6.5. We see that the amplitude of the oscillation is maximum
when the external frequency ω = ω0 , and that there is a sharp change in the amplitude and
in the phase at this frequency. This is our first demonstration of amplitude resonance, where
there is a large response from the oscillator when it is driven at a frequency that matches
the natural frequency of the oscillator. We will comment on resonance in more detail later
in this section.
20 5
4
Phase difference
15
Amplitude A
3
10
2
5
1
0 0
0.0 0.5 1.0 1.5 2.0 0.0 0.5 1.0 1.5 2.0
External frequency ω/ωo External frequency ω ωo
Figure 6.5: The amplitude A and phase difference φ as a function of the external frequency
ω, in the special case of driven simple harmonic motion with no damping.
A more realistic situation is the case where damping is not zero. For this case, the
standard method is to try substituting a solution x = A cos(ωt − φ) in (6.6.1). The choice
of a cosine solution is motivated by the fact that the right-hand side of (6.6.1) contains a
cosine term with a frequency ω. We need a function x(t) whose derivatives will ultimately
give a cosine with the same frequency. We know that the first derivative of the cosine is
sine. Therefore we introduce a phase term φ, which if chosen properly will allow all of the
terms on the left-hand side of (6.6.1) to add up to a cosine, and match the right-hand
side. To simplify the mathematics, we write our trial solution as a complex exponential,
x = A ei(ωt−φ) , and also replace the term D cos(ωt) in this equation with D eiωt . We will
carry out the algebra
and then at the end, we will take the real part of the solution, because
cos(ωt) = Re eiωt . By substituting and taking the derivatives, we obtain:
i i
i i
i i
By canceling out the eiωt from all terms and multiplying by eiφ = cos φ + i sin φ on both
sides, we obtain:
− Aω 2 + 2γωiA + ω02 A = D eiφ = D (cos φ + i sin φ) (6.6.11)
By equating the real parts on the two sides of this equation and also equating the imaginary
parts on the two sides to each other, we obtain:
A ω02 − ω 2 = D cos φ (6.6.12)
dx dxp
v= = − Aω sin (ωt − φ) (6.6.15)
dt dt
F0 /m
A= q 2 (6.6.16)
ω02 − ω 2 + 4γ 2 ω 2
2γω
tan φ = (6.6.17)
ω02 − ω 2
The Python code in Example 6.5 numerically solves (6.6.1) using the Euler method
described in Chapter 2. The results are plotted in Figure 6.6, where the left graph is a
plot of x(t) and the right graph is a plot of v(t). Notice that in Figure 6.6, the transient
behavior is dominant for t < 60 s, and that the sinusoidal steady-state behavior dominates
for t > 60 s.
Example 6.5: Python code for driven harmonic oscillator
Write a Python code which numerically solves the differential equation for the driven
SHO, ẍ + 2γ ẋ + ω02 x = D cos(ωt), by using the Euler method described in Chapter 2. Use
the numerical values m = 1 kg, k = 1 N/m , b = 15 Ns/m , D = 5 m/s2 and ω = 3 s−1 , from
t = 0 to t = 80 s. Plot both the position x(t) and the velocity v(t). Discuss the behavior of
the solution at small times and at large times. Use initial conditions x(0) = 2 and v(0) = 0.
Solution:
The following Python code clearly shows the change of behavior of the solutions x(t)
and v(t) as time progresses. The initial part of the functions contains an exponentially
decaying transient component. At longer times, the solutions approach the steady-state
solutions, which oscillate with time, and with a constant amplitude. Note that the red (in
the e-book) curled arrows in the code denote a line continuation, i.e., the text following a
curly arrow is a continuation of the line above it. Those arrows are not to be included in
your code.
i i
i i
i i
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
#d e f i n e i n i t i a l c o n d i t i o n s and a r r a y s
x = np . z e r o s ( num steps )
v = np . z e r o s ( num steps )
x [ 0 ] = 2.0
v [ 0 ] = 0.0
#perform t h e E u l e r s t e p a l g o r i t h m
f o r i in range ( 1 , num steps ) :
x [ i ] = x [ i −1] + dt ∗ v [ i −1]
v [ i ] = v [ i −1] + dt ∗(−k/m ∗x [ i ]− dt ∗b/m ∗v [ i −1]+5∗ np . c o s
,→ ( 3 ∗ dt ∗ i ) )
time = np . l i n s p a c e ( 0 , tmax , num steps )
f i g , a x e s = p l t . s u b p l o t s ( 1 , 2 , f i g s i z e =(10 ,4) )
a x e s [ 0 ] . p l o t ( time , x )
a x e s [ 0 ] . s e t x l a b e l ( ” time , s ” , f o n t s i z e =18)
a x e s [ 0 ] . s e t y l a b e l ( ”x , m” , f o n t s i z e =18)
a x e s [ 1 ] . p l o t ( time , v )
a x e s [ 1 ] . s e t x l a b e l ( ” time , s ” , f o n t s i z e =18)
a x e s [ 1 ] . s e t y l a b e l ( ”v , m/ s ” , f o n t s i z e =18)
f i g . s a v e f i g ( ” ex4 . 5 . e p s ” )
4
2 3
2
1
1
v, m/s
x, m
0
0
−1
−2
−1
−3
−4
−2
0 20 40 60 80 0 20 40 60 80
time, s time, s
Figure 6.6: The result of running the Python code in Example 6.5 for the externally driven
harmonic oscillator. The effect of the transient part of the solution is clearly seen for times
t < 60 s, while the steady-state solution is dominant for t > 60 s.
i i
i i
i i
In addition, we can plot the amplitude A and phase angle φ as a function of frequency
using (6.6.16) and (6.6.17). The plots of the amplitude A and of the phase angle φ as a
function of the external frequency ω are shown in Figure 6.7. When ω << ω0 , we find that
the phase difference φ goes to zero, while for ω >> ω0 we find that the phase goes to π (180
degrees) .
Figure 6.7: Plots of the amplitude A (ω) and phase angle φ (ω) in the case of externally-
driven damped harmonic motion.
One of the most striking results shown in Figure 6.7 is the amplitude A (ω) as a function
of frequency. Notice that similar to the undamped case, amplitude resonance can occur.
The damping prevents the amplitude from blowing up to infinity at resonance. Although
the small value of γ = 0.01 results in a large amplitude at resonance, it is not infinite; it
simply is beyond the scale of the graph here. Notice that the amplitude resonance occurs
when ω ≈ ω0 , because of the small values of γ used to make Figure 6.7. Next, we will
compute the exact frequency at which resonance occurs.
The resonance frequency is the value of the drive frequency ω, which causes amplitude
resonance, and is found by computing:
dA
=0 (6.6.18)
dω
ω=ωr
q
We find that A is a maximum when ω = ωr = ω02 − 2γ 2 , where ωr is the resonant frequency.
The maximum value of A defined as Amax is then:
F /m
Amax = q0 (6.6.20)
2γ ω02 − γ 2
i i
i i
i i
OUTPUT:√ √
Fo((−ω 2 +ωo2 )Cos[tω]+2γωSin[tω])
−t γ+ γ 2 −ωo2 t −γ+ γ 2 −ωo2
e C[1] + e C[2] +
m 4γ 2 ω 2 +(ω 2 −ωo2 )
2
steady = x/.Exp[t ] → 0
trans = x − steady
√ √
−t γ+ γ 2 −ωo2 t −γ+ γ 2 −ωo2
OUTPUT: e C[1] + e C[2]
m = 1; γ = 0.1; Fo = 1; ωo = 1;
u = Table[{ω, NMaxValue[steady, t]}, {ω, 0, 3, .02}];
ListPlot[u, FrameLabel → {“Driving frequency ω, Hz”, “Amplitude A(ω)”},
Frame → True, BaseStyle → (FontSize → 18), Joined → True, ImageSize → Large]
5
4
Amplitude A(ω)
0
0.0 0.5 1.0 1.5 2.0 2.5 3.0
Driving frequency ω, Hz
i i
i i
i i
1 (F0 /m)2
<E>= mω 2 + mω02 (6.7.3)
4
2
ω02 − ω 2 + 4γ 2 ω 2
and finally:
F02 ω 2 + ω02
<E>= (6.7.4)
4m ω 2 − ω 2 2 + 4γ 2 ω 2
0
A plot of < E > as a function of the external driving frequency ω is shown in Figure
6.8. Notice that there is a peak similar to the peak we saw for the amplitude graph A(ω)
in Figure 6.7.
Figure 6.8 is illustrating energy resonance, and the energy resonance frequency is the
drive frequency which maximizes the average energy of the oscillator.
In the context of energy resonance and driven oscillations, the Q factor is defined as the
“frequency-to-bandwidth ratio” of the resonator, thus:
ωr
Q≡ (6.7.5)
∆ω
where ωr is the resonant angular frequency of the peaked graph in Figure 6.7. ∆ω in this
definition is the full width at half maximum (FWHM) of the energy resonance peak. In
i i
i i
i i
order to find ∆ω, we need to find the two points ω1 and ω2 on this graph, which correspond
to half of the maximum of the peak.
Figure 6.8: Plot of the time averaged energy < E > from (6.7.4), as a function of the
driving frequency ω. The value of γ = 0.03 results in a lightly damped harmonic motion,
corresponding to a large Q-factor.
To calculate Q, we assume a system with small damping (i.e., a sharp resonance), and
2
we substitute ω ' ω0 in (6.7.4), except for the term ω02 − ω 2 , to obtain:
F02 ω02 1
<E>= (6.7.6)
2m ω 2 − ω 2 2 + 4γ 2 ω 2
0 0
By approximating ω0 ' ω, the term ω02 − ω 2 = (ω0 + ω) (ω0 − ω) ' 2ω0 (ω0 − ω), and the
previous equation becomes:
F02 1
<E>= (6.7.7)
8m (ω0 − ω)2 + γ 2
Next, taking the derivative of < E > with respect to the driving frequency ω and setting
it equal to zero, we find that the maximum of this function occurs at ω = ω0 . Therefore,
the average energy resonance occurs when ω = ω0 . By substituting this value in (6.7.4), the
maximum time average of the total energy is given by:
F02
< E >max = (6.7.8)
8mγ 2
Next, we need to find the frequency at which the average energy < E > drops to half its
maximum value. We substitute (6.7.8) in (6.7.7) to obtain:
1 1 F02 F02 1
< E >max = = (6.7.9)
2 2 8mγ 2 8m (ω0 − ω)2 + γ 2
ω1 = γ + ω0 ω2 = ω0 − γ (6.7.11)
i i
i i
i i
∆ω = ω1 − ω2 = 2γ (6.7.12)
Therefore from the definition of the Q-factor in (6.7.5), we find the ratio of the resonance
frequency ω0 over the FWHM ∆ω:
This is the same simple equation we derived for underdamped motion, and involves only
the coefficients of the second-order differential equation describing most resonant systems,
electrical or mechanical. In electrical systems, the stored energy is the sum of energies stored
in lossless inductors and capacitors; the lost energy is the sum of the energies dissipated
in resistors per cycle. In mechanical systems, the stored energy is the maximum possible
stored total energy, i.e., the sum of the potential and kinetic energies at some point in time,
while the lost energy is that lost to damping forces such as friction.
i i
i i
i i
sol = DSolve[{x00 [t] + ωo∧ 2 ∗ x[t] + 2 ∗ γ ∗ x0 [t] == (Fo*/m ∗ Cos[ω ∗ t]}, x[t], t];
x = x[t]/.sol[[1]]//Simplify
OUTPUT:√ √
Fo((−ω 2 +ωo2 )Cos[tω]+2γωSin[tω])
−t γ+ γ 2 −ωo2 t −γ+ γ 2 −ωo2
e C[1] + e C[2] +
m 4γ 2 ω 2 +(ω 2 −ωo2 )
2
steady = x/.Exp[t ] → 0
en = (1/2) ∗ m ∗ steady∧ 2
m = 1; γ = 0.1; Fo = 1; ωo = 1;
u = Table[{ω, NMaxValue[en, t]}, {ω, 0, 3, 0.02}];
ListPlot[u, PlotRange → All, FrameLabel → {“Driving frequency ω, Hz”, “Energy E(ω)”},
Frame->True, LabelStyle → {FontSize → 20}, Joined → True, ImageSize → Medium]
12
10
Energy E(ω)
8
6
Out[ ]=
4
2
0
0.0 0.5 1.0 1.5 2.0 2.5 3.0
Driving frequency ω, Hz
Figure 6.9: RLC series circuit with external voltage V (t), resistance R, inductance L, and
capacitance C.
i i
i i
i i
The governing differential equation can be found by substituting into Kirchhoff’s voltage
law (KVL), the constitutive equation for each of the three elements. From the KVL,
V (t) + VR + VL + VC = 0 (6.8.1)
where VR = −IR is the voltage drop across the resistor, VL = is the voltage drop
−L dI
dt
Q
across the inductor, VC = − C is the voltage drop across the capacitor, and V (t) is the
time varying voltage from the source. Note that we are using the normal conventions that
I is the current in the circuit and Q is the charge on the capacitor. Furthermore, L is the
inductance of the inductor, R is the resistance of the resistor, C is the capacitance of the
capacitor, and the current is I = Q̇. For simplicity, we will assume that V (t) = V0 cos(ωt).
Substituting these expressions in the KVL equation:
dI Q
L + IR + = V0 cos(ωt) (6.8.2)
dt C
Dividing by L leads to the second-order differential equation:
General Differential Equation for RLC Series Circuit
d2 Q(t) R dQ(t) 1 V0
2
+ + Q(t) = cos(ωt) (6.8.3)
dt L dt LC L
This equation has the exact same mathematical form as equation (6.6.1) for a driven
damped harmonic oscillator:
d2 x b dx k F0
+ + x= cos(ωt) (6.8.4)
dt2 m dt m m
where (6.6.1) is written without the substitution γ = b/2m.
Because (6.8.3) and (6.6.1) have the same mathematical form, we can reuse the solution
of (6.6.1) in order to get the solution of (6.8.3). The mathematical similarities imply that the
physical behaviors of the two systems are similar. The forced damped harmonic oscillator
has two solutions: a transient solution, which decays as t → ∞, and a steady-state solution
consisting of a sinusoidal oscillation with a frequency equal to the drive frequency. Physically,
these solutions are describing the displacement of the mass, m. Because of the similarities
between (6.8.3) and (6.6.1), we know that the RLC circuit will also have two solutions, a
transient, which will also decay, and a steady-state sinusoidal solution, which has the same
frequency as the voltage source in the circuit. In the case of the RLC circuit, it is the value
of the charge on the capacitor that is oscillating.
Furthermore, we can then draw exact analogies between mechanical systems and elec-
trical systems. These analogies are shown in Table 6.1 and are found by comparing the
locations of the variables in (6.8.3) to those in (6.8.4). The simplest analogy is of that
between the displacement x of the mass and the charge Q on the capacitor. A more inter-
esting analogy exists between L and m. The inductance L in (6.8.3) appears in the same
location as m in (6.6.1). This implies that L is taking on a role similar to that of mass. This
analogy makes physical sense if you recall that the inductance L measures the inductor’s
ability to resist changes in current (Q̇), while the mass m is a measure of inertia, an object’s
ability to resist changes in velocity (ẋ).
It is also possible to use the analogies to identify qualitative characteristics of the RLC
circuit’s behavior. For example, we know that the coefficient of x in (6.6.1) is the angular
frequency ω02 of oscillation for an undamped oscillator. We can conclude from (6.8.3) that
the oscillation frequency of a LC circuit (where R = 0) will be given by:
i i
i i
i i
Mass m Inductance L
Spring Constant k Inverse Capacitance 1/C
Damping γ Resistance R
q q
Natural Frequency k
m
1
LC
Table 6.1: Analogy between electrical RLC circuit in series and the mechanical driven har-
monic oscillator.
Finally, we can derive an expression about the energy of the RLC circuit. We begin by
multiplying (6.8.3) by Q̇ and L:
1
LQ̇Q̈ + RQ̇2 + QQ̇ = V0 cos(ωt)Q̇ (6.8.6)
C
Notice that:
1 d 1 d
Q̇Q̈ = Q̇2 and QQ̇ = Q2
2 dt 2 dt
Substitution of the previous equations into (6.8.6) gives:
d 1 11 2
LQ̇2 + Q = −RQ̇2 + V0 cos(ωt)Q̇ (6.8.7)
dt 2 2C
The 12 LQ̇2 and 12 C1 Q2 terms in (6.8.7) are the energy stored in the inductor and capacitor,
respectively. We can see that if R = 0 and if there is no voltage source in the circuit V0 = 0,
then the right-hand side of (6.8.7) is zero, so the total energy of the circuit is conserved.
The terms on the right-hand side describe how energy is either lost or gained in the system.
Recall from your introductory physics course that the power dissipated by a resistor is
P = −I 2 R, which is the first term on the right-hand side of (6.8.7). The second term on the
right-hand side of (6.8.7) is describing the rate at which energy is supplied by the variable
voltage source. Notice that V0 cos(ωt)Q̇ is in the form of V I, which is the formula for electric
power.
Analogies can be a powerful tool in physics, which can be used to take an understanding
of one system, and apply it to another mathematically similar system. As with any analogy
(physics-related or not), be careful not to take the analogy too far. Doing so can lead to
confusion at best and to a complete misunderstanding of the system at worst.
i i
i i
i i
however, such a definition could be difficult to work with analytically. We should note
that modern CAS like Mathematica can handle such equations in a variety of applications,
however, there is much value in learning how to work with such cases analytically. So, the
question remains: how do we represent such periodic discontinuous functions in a way that
lends itself to analytical solutions? The answer, which we will demonstrate in this section,
is that we can represent periodic discontinuous functions as a sum of sines and cosines
called a Fourier series. This interesting result is known as Fourier’s theorem and applies
to periodic functions that are continuous, or have a finite number of discontinuities over
the period τ . Since most functions describing physical processes have only a finite number
of discontinuities over an interval t ∈ [0, τ ], we can consider Fourier’s theorem to be widely
applicable.
Fourier’s theorem is only one part of the solution of handling discontinuous periodic
drive functions. By representing drive functions as a series of sines and cosines, we need to
be able to solve equations of the form:
In other words, we need to know how to solve linear ordinary differential equations that
have multiple terms on the right-hand side. Because the differential equation is linear in x
and its derivatives, we can use the powerful concept of the principle of superposition, which
will allow us to break up the differential equation into pieces and then sum the solutions of
the individual pieces.
d2 x
2
dx d d
a 2 + b + cx = a 2 + b + c x = F (t) (6.9.3)
dt dt dt dt
which is very similar to the equations describing oscillations that we have worked with so
far. We can rewrite such an equation as L̂(x) = F (t), where the differential operator L̂ is
defined as:
d2 d
L̂ ≡ a 2 + b + c (6.9.4)
dt dt
The operator L̂ is linear if it satisfies:
where c1 and c2 are arbitrary constants, and x1 and x2 are two solutions of the differential
equation. It is straightforward to show that L̂ defined in (6.9.4) is a linear operator.
i i
i i
i i
where L̂ is defined as in (6.9.4). Suppose we know two functions x1 and x2 , such that
L̂ (x1 ) = F1 (t) and L̂ (x2 ) = F2 (t). In other words, x1 is the solution to (6.9.6) without F2 (t)
on the right-hand side, and x2 is the solution without F1 (t) present. Then:
Therefore, the sum of the individual solutions x1 and x2 is the solution to the full differential
equation (6.9.6). We can, of course, extend this to an arbitrary number of terms on the
right-hand side:
N
X
L̂ (x) = ci Fi (t) (6.9.8)
i=1
where xi satisfied L̂ (xi ) = Fi (t) and xc is the complementary solution of the ODE. The
ability to add solutions in the above manner is called the principle of superposition and
applies only to linear differential equations. In Chapter 13, we will study nonlinear systems,
and the principle of superposition will not apply to those systems.
The principle of superposition will be critical in dealing with Fourier series representa-
tions of discontinuous drive forces. When we represent the discontinuous drive force as a
Fourier series, the resulting differential equation will be of the form:
We now know that in order to solve this type of equation, all we need to do is solve the
system of equations:
for i = 0, . . . , N . Note that for i = 0, cos(0) = 1, and we obtain the solution for the constant
term on the right-hand side of the differential equation. Similarly sin(0) = 0 gives the homo-
geneous solution. While there may be a large number of equations to solve, there are ways
that the process can be generalized. For example, notice that in the above case, by changing
the value of i, we are simply changing the drive frequency in the case of i > 0. Next, we
need to learn how to represent a discontinuous periodic function as a Fourier series.
i i
i i
i i
where the lowest angular frequency is ω = 2π/τ , and the higher frequencies are integer
multiples of the fundamental angular frequency ω. The coefficients an and bn represent
the amplitudes of the various waves, and are known as Fourier coefficients. The Fourier
coefficients an and bn are computed by using the following general integral relationships:
(
0 m 6= n
Z τ /2
cos (nωt) cos (mωt) dt =
−τ /2 τ /2 m = n 6= 0
(
0 m 6= n
Z τ /2
sin (nωt) sin (mωt) dt = (6.9.13)
−τ /2 τ /2 m = n 6= 0
Z τ /2
cos (nωt) sin (mωt) dt =0 for all integers n and m
−τ /2
The calculation involves multiplying both sides of (6.9.12) with either cos(2πmt/τ ) or
sin(2πmt/τ ), and integrating over the whole period. The results of this calculation are
known as Fourier’s trick to many generations of physicists:
2 τ /2
Z
an = f (t) cos(nωt)dt (6.9.14)
τ −τ /2
2 τ /2
Z
bn = f (t) sin(nωt)dt (6.9.15)
τ −τ /2
It is not always possible to compute coefficients of all values of n, so a finite Fourier series
approximates the original function f (t), and the approximation improves as N increases.
The Fourier series is generally presumed to converge everywhere except at discontinuities,
where it converges to the mean of the values of f (t) before and after the discontinuity.
Solution:
Before finding the Fourier series, it is often helpful to plot the function f (x).
i i
i i
i i
1.0
0.5
f(x)
0.0
-0.5
-1.0
-5 0 5
Position x
Although in this case the independent variable is x, the function f (x) is periodic in x, so
all of the formulas for Fourier series still apply. Next, it is important to identify the period.
As can be seen from the graph (and the definition of f ), the period is τ = 2π. Using the
formulas (6.9.14) and (6.9.15), we can compute the Fourier coefficients with ω = 2π/τ = 1.
First we find the an coefficients:
1 π x
Z
an = cos( 2πnx
2π ) dx
π −π π
2 Z π
1
= x cos(nx) dx
π −π
2 π
1 x sin(nx)
=
π n
−π
therefore we see that an = 0 for all values of n, because sin(nπ) = 0 for all integers n. Next,
we compute the bn coefficients.
1 π x
Z
bn = sin( 2πnx
2π ) dx
π −π π
2 Z π
1
= x sin(nx) dx
π −π
2 2
=− cos(nπ) + 2 2 sin(nπ)
πn π n
Again, we use sin(nπ) = 0 for integer values of n. However, note that cos(nπ) = (−1)n .
Therefore, we find that:
2 (−1)n+1
bn =
πn
The resulting Fourier series is therefore:
∞
2 X (−1)n+1
f (x) = sin(nx)
π n
n=1
When x = ±π, the Fourier series converges to 0, which is the mean of the left and right
limits of f (x) at x = ±π.
We can use Python to numerically compute the Fourier coefficients. Using a computer
to numerically find Fourier coefficients can be helpful in cases where the function is compli-
i i
i i
i i
0.5
0.0
−0.5
−1.0
−3 −2 −1 0 1 2 3
x
Figure 6.10: Results from Python code for calculating the Fourier coefficients an , bn of a
periodic function. The graph shows a comparison of the original function (orange in the
e-book) and the approximation using the Fourier series (blue in the e-book).
cated. With a computer, we can find as many terms in the Fourier series as we need (and
have the resources/time for!).
Algorithm 5 uses SciPy’s quad command to find the first 10 terms of the Fourier series
of the sawtooth function from Example 6.8. Note that quad returns two arguments, the first
being the value of the integral and the second argument being the error. Notice that we
defined f (x) as its own function. In addition, we also defined two functions a coefficient and
b coefficient, which are functions of n, and correspond to an and bn . The coefficients an and
bn are numbers found by integrating f (x) cos(2πnx/tau) and f (x) sin(2πnx/tau), respec-
tively. In order to perform the integrals over x, quad needs the integrand to be defined as a
function of x, which is why we defined the additional functions a integrand and b integrand.
We then store the Fourier coefficients as an array. Finally, we need to plot the result-
ing Fourier series. In order to make the plot, we need to evaluate the Fourier series at
several values of x. In Algorithm 5, we used 100 values of x between x = −tau/2 = −π
and x = tau/2 = π. Recall that in NumPy if we compute a function of an array, such as
cos ([1, 2, 3]) we obtain [cos(1), cos(2), cos(3)], hence the variable fourier series is actually an
array. The definition of fourier series is a bit complex, so make sure you understand what
is being done. Finally, we plot the Fourier series and the original function together, and the
result is shown in Figure 6.10, where f (x) is the orange line and the 10-term Fourier series
is in blue. If we increased the number of terms (by changing the value of num fc), we would
get a better approximation.
i i
i i
i i
import numpy a s np
from s c i p y . i n t e g r a t e import quad
import m a t p l o t l i b . p y p l o t a s p l t
tau = 2 . 0 ∗ np . p i #p e r i o d o f f u n c t i o n
num fc = 10 #number o f terms i n F o u r i e r s e r i e s
def f ( x ) : #f u n c t i o n f o r which t h e F o u r i e r s e r i e s i s c a l c u l a t e d
return x/np . p i
def a c o e f f i c i e n t ( n ) : #F o u r i e r c o e f f i c e n t , a n
def a i n t e g r a n d ( x ) : return f ( x ) ∗np . c o s ( 2 . 0 ∗ np . p i ∗n∗x/ tau )
i n t v a l u e , e r r o r = quad ( a i n t e g r a n d ,− tau / 2 . 0 , tau / 2 . 0 )
return 2 . 0 / tau ∗ i n t v a l u e
def b c o e f f i c i e n t ( n ) : #F o u r i e r c o e f f i c e n t , a n
def b i n t e g r a n d ( x ) : return f ( x ) ∗np . s i n ( 2 . 0 ∗ np . p i ∗n∗x/ tau )
i n t v a l u e , e r r o r = quad ( b i n t e g r a n d ,− tau / 2 . 0 , tau / 2 . 0 )
return 2 . 0 / tau ∗ i n t v a l u e
#compute a n ’ s and b n ’ s
a = np . a r r a y ( [ a c o e f f i c i e n t ( n ) f o r n in range ( 0 , num fc ) ] )
b = np . a r r a y ( [ b c o e f f i c i e n t ( n ) f o r n in range ( 0 , num fc ) ] )
#compute t h e F o u r i e r s e r i e s as a f u n c t i o n o f x f o r one p e r i o d
num xvals = 100 #number o f v a l u e s
x v a l s = np . l i n s p a c e (−tau / 2 . 0 , tau / 2 . 0 , num xvals )
f o u r i e r s e r i e s = a [ 0 ] / 2 . 0 + sum ( [ a [ n ] ∗ np . c o s ( 2 . 0 ∗ np . p i ∗n∗ x v a l s /
tau ) + b [ n ] ∗ np . s i n ( 2 . 0 ∗ np . p i ∗n∗ x v a l s / tau ) f o r n in range ( 1 ,
num fc ) ] )
p l t . show ( )
Algorithm 5: Python algorithm for evaluating Fourier coefficients for the sawtooth
function.
i i
i i
i i
have m as a term in the formula of the driving force, but it will simplify the mathematics
later.
The driving force in (6.9.16) acts on a damped harmonic oscillator with damping param-
eter γ and natural frequency ω0 , whose equation of motion is then:
F (t)
ẍ + 2γ ẋ + ω02 x = (6.9.17)
m
In order to solve for x(t) using the analytical methods described in this chapter, we will
need to use the Fourier series for F (t). From Example 6.8, we know that:
∞
2mA X (−1)n+1
F (t) = sin(nωt) (6.9.18)
π n
n=1
where we changed the independent variable from x to t, included the multiplicative constant
A, and explicitly included ω into the Fourier series. Note that ω = 1 because τ = 2π. So the
resulting differential equation to solve is:
∞
2A X (−1)n+1
ẍ + 2γ ẋ + ω02 x = sin(nωt) (6.9.19)
π n
n=1
However, the work done in this chapter had cosine terms on the right-hand side of the
differential equation. In order to use the solutions developed in this chapter, we will need
to use the identity: sin(a) = cos(a − π/2). Doing so results in (6.9.19) becoming:
∞
2A X (−1)n+1
ẍ + 2γ ẋ + ω02 x = cos(nωt − π/2) (6.9.20)
π n
n=1
2A (−1)n+1
ẍn + 2γ ẋn + ω02 xn = cos(nωt − π/2) (6.9.21)
π n
We use trial solutions of the form xn (t) = An cos(nωt − π/2 − φn ) and follow the procedure
outlined for the solution of the driven damped harmonic oscillator. The result is:
2A(−1)n+1
An = q 2 (6.9.22)
nπ ω02 − (nω)2 + 4γ 2 (nω)2
2γnω
tan φn = (6.9.23)
ω02 − (nω)2
Notice that the above solution is the same as (6.6.16) and (6.6.17), however the constant
coefficient is now 2A(−1)n+1 /nπ, and the drive frequency is nω instead of ω. Finally, we
can use the principle of superposition to find the solution to (6.9.20):
∞
2A X (−1)n+1
x(t) = q cos(nωt − π/2 − φn ) (6.9.24)
nπ 2
n=1 ω02 − (nω)2 + 4γ 2 (nω)2
where φn is defined in (6.9.23).
i i
i i
i i
x2 v2
+ =1 (6.10.4)
2E/k 2E/m
and the semimajor y-axis equal to b = 2E/m. Examples of harmonic oscillators with
different energies, E, are shown in Figure 6.11. In this case, we see that ellipses with larger
area correspond to the trajectories with larger energies.
1
Figure 6.11: Phase space of simple harmonic oscillator.
Phase space diagrams can tell us something about the motion of the system. The curve
in a phase diagram is sometimes referred to as a trajectory in phase space. An image of all
i i
i i
i i
possible trajectories in a phase space is sometimes called a phase portrait or phase diagram.
Each point along a trajectory gives the system’s position and velocity. The presence of
closed-loop trajectories in a phase space diagram is evidence of oscillatory motion. A closed-
loop trajectory tells us that after some time τ , the system returns to its initial position and
velocity. The elliptical trajectories in Figure 6.11 tell us that the SHO displays oscillatory
motion (which we already knew!).
Next, we will look at the phase space diagram of an underdamped harmonic oscillator.
In this case, we write:
x(t) = Ae−γt cos(ωd t − φ) (6.10.5)
v(t) = dx/dt = −Aωd e−γt sin(ωd t − φ) − Aγe−γt cos(ωd t − φ) (6.10.6)
It can be difficult to visualize the trajectories directly from the above equations. Therefore,
we can make the following change of coordinates, in order to simplify the algebra:
u = ωd x, w = γx + v (6.10.7)
i i
i i
i i
A = 1; ωd = 1; φ = 0; γ = 0.1;
x = A ∗ Exp[−γ ∗ t]Cos[ωd ∗ t − φ];
v = D[x, t];
0.5
Out[ ]= x
-0.5 0.5 1.0
-0.5
The above algorithm can be used to generate the phase diagram for any system whose
position x(t) is known. Recall that the D in the above code is Mathematica’s command for
differentiation. While the phase space diagram does not provide new information about the
behavior of the harmonic oscillator, we will find in Chapter 13 that phase space diagrams
will become a critical tool for understanding the possible behaviors of a nonlinear system.
In Chapter 13, we will develop a method of finding a system’s phase space diagram, without
solving analytically the system’s equations of motion. It is often the case in nonlinear
systems that the equations of motion cannot be solved analytically. Therefore the phase
space diagram may be the only tool available in getting information about the behavior of
the system.
i i
i i
i i
d2 x dx k
q
−γt −ωt
+ 2γ + x=0 x(t) = e ωt
+ A1 e ω = γ 2 − ω02
A0 e
dt2 dt m
where the constant γ = b/2m is the damping
q parameter. There are three possible types of
solutions, depending on whether ω = γ − ω02 is real, imaginary, or zero: (a) overdamped
2
oscillations (γ > ω0 ), (b) underdamped oscillations (γ < ω0 ), and (c) critically damped
oscillations (γ = ω0 or ω = 0). The corresponding solutions are:
q
x(t) = A0 eωt + A1 e−ωt e−γt ω= γ 2 − ω02 γ > ω0 Overdamped
q
x(t) = Ae−γt cos (ωd t − φ) ωd = ω02 − γ 2 γ < ω0 Underdamped
−ω0 t
x(t) = (A + Bt) e γ = ω0 Critically Damped
The energy and energy loss for a weakly damped oscillator are:
1 dE
E = kA20 e−2γt γ << ωd = −2γE (6.11.1)
2 dt
The dimensionless Q-factor expresses the degree of damping present in an oscillator, with
higher Q indicating a lower rate of energy loss:
ω0
Q=
2γ
The forced damped oscillator is described by the amplitude A (ω) and phase angle φ (ω),
both are functions of the external driving frequency ω:
F0 /m 2γω
A= q tan φ =
2 ω02 − ω 2
ω02 − ω 2 + 4γ 2 ω 2
where xc (t) is one of the three possible solutions of the homogeneous equation. This
oscillator exhibits the phenomenon of resonance, with the following resonance frequency
and resonance amplitude:
F /m
q q
ωr = ωd2 − γ 2 = ω02 − 2γ 2 Amax = q0
2γ ω02 − γ 2
i i
i i
i i
There is an analogy between the charge Q(t) in electrical RLC circuits and the position
x(t) in mechanical oscillators:
d2 Q(t) R dQ(t) 1 d2 x b dx k
+ + Q(t) = 0 + + x=0
dt2 L dt LC dt2 m dt m
The general mathematical form of a Fourier series for any periodic function is:
N
a0 X
f (t) = + an cos(nωt) + bn sin(nωt)
2
n=1
2 τ /2
2 τ /2
Z Z
an = f (t) cos(nωt)dt bn = f (t) sin(nωt)dt
τ −τ /2 τ −τ /2
Phase space diagrams are important in analyzing nonlinear systems. They are plots of
the position x(t) of a particle on the x-axis and its corresponding speed v(t) (or
momentum) on the y-axis.
2. Starting from the general expression x(t) = A cos ( ω0 t) + B sin ( ω0 t) for a harmonic
oscillation, show that an equivalent form
√ of this equation is given by the expressions:
x(t) = C cos ( ω0 t + φ) where C = A2 + B 2 tan φ = B/A.
3. A mass m is attached to a spring of spring constant k and moves with underdamped
harmonic motion along the x-axis, with the air resistance given by F = −bv.
(a) Show that the difference in times between two successive maxima is given by:
2π
∆t = p
ωo2 − γ 2
(b) Show that the ratio of the amplitudes of two successive maxima of the position is
given by:
R = eγ∆t
i i
i i
i i
mass, where ω = k/m. Describe qualitatively the position and the speed of this
p
mass (x(t), v(t)) as a function of time t. You may neglect any damping present in the
system.
5. A spring has a spring constant of k = 8 N/m and a mass m = 0.5 kg is attached to
it. At time t = 0 a force F = 6 sin (4t), in Newtons, is applied to the mass. Find an
analytical expression for the position x (t) at time t.
6. The differential equation for an oscillating system is:
d2 z
− 16z − 160 cos (6t) = 0
dt2
with the initial conditions z(0) = 0 and dt (0)
dz
= 0. Plot the solution z(t) for several
periods and interpret the results.
7. A cylinder of mass m and radius R floats with its axis vertically in a liquid of density
ρ. The cylinder is given a small vertical displacement downwards and is released. Find
the period of oscillation.
8. A mass m moves in a region of xy-plane where the force is given by F = −kxî − ky ĵ.
(a) Show that depending on the initial conditions of the problem, the motion will be
an ellipse, a parabola, or a hyperbola.
(b) Find the potential energy of this oscillator.
9. Consider a mass m moving on a vertical wire inside a gravitational field. The shape of
the wire is a cycloid, described by the equations: x = aθ − a sin θ and y = −a + a cos θ.
(a) Find the speed of the mass at the bottom of the cycloid wire, if it is released at
the origin (0, 0) at time t = 0.
(b) Find the period of oscillation.
10. A mass m which is on a frictionless table is connected to two fixed points on opposing
walls, by two springs of equal natural length, of negligible mass and spring constants
k1 and k2 respectively. The mass is displaced horizontally and then released. Prove
that the period of oscillation is τ = 2π m/(k1 + k2 )
p
11. A spring having a spring constant k and negligible mass has one end fixed on an
inclined plane of angle θ and a mass m at the other end, as shown in Figure 6.12. If
the mass m is pulled down a distance x below the equilibrium position and released,
find the displacement from the equilibrium position at any time if (a) the incline is
frictionless, (b) the incline has coefficient of friction µ.
i i
i i
i i
12. A particle moves with simple harmonic motion along the x-axis. At times t0 , 2t0 , 3t0 ,
it is located at x
= a, b, c, respectively. Prove that the period of oscillation is T =
2πt0 / cos−1 a+c .
2b
13. Two equal masses m are connected by springs having equal spring constant k, as shown
in Figure 6.13, so that the masses are free to slide on a frictionless. The walls attached
to the ends of the springs are fixed. Set up the differential equations of motion of the
masses.
(a) Find the possible frequencies of oscillation, also called the normal frequencies.
(b) Suppose that the first mass is held at its equilibrium position while the second
mass is given a displacement of magnitude a > 0 to the right of its equilibrium
position. The masses are then released. Find the position of each mass at any
later time.
Figure 6.13: Problem 6.13: Two coupled massess connected to fixed ends.
14. Two equal masses m on a horizontal frictionless table, as shown in the Figure 6.14, are
connected by springs of equal spring constants k. The end of one spring is attached
to a wall and the masses are set into motion.
(c) Set up the equations of motion of the system.
(d) Find the possible frequencies of oscillation, also called the normal frequencies.
15. Two masses, m1 and m2 , are attached to the ends of a spring of constant k, which
is on a horizontal frictionless table as shown in Figure 6.15. If the masses are pulled
apart and thenpreleased, prove that they will vibrate with respect to each other with
period T = 2π µ/k where µ = M1 M2 /(M1 + M2 ) is the reduced mass of the system.
Figure 6.15: Problem 6.15: Two masses attached to a spring on a frictionless table.
i i
i i
i i
16. Find the frequencies of vibration for the system of particles of equal masses m con-
nected by springs with the same spring constant k, as indicated in Figure 6.16.
Figure 6.16: Problem 6.16: Two coupled masses hanging from verticle springs.
17. A pendulum has length L = 1 m, mass m = 0.1 kg and the gravitational acceleration
is g = 9.8 m/s2 .
(a) Find the kinetic energy T and potential energy V as a function of time t, when
the pendulum is released from rest at an angle of 3 degrees from equilibrium. Plot
T, V and the sum T + V on the same graph using the programming language of
your choice.
(b) Repeat part (a) when the pendulum is released from rest at an angle of 80 degrees
from equilibrium. How are the graphs of T, V , T + V different from part (a)?
(c) Find the time averages < T > and < V > over a complete period of oscillation
in part (a) and compare them with the total mechanical energy E = T + V .
18. A pendulum has length L = 1 m, mass m = 0.1 kg, and is released from rest at an
angle of 3 degrees from equilibrium. Use a value of g = 9.8 m/s2 for the gravitational
acceleration.
(a) Find the tension T of the string as a function of time t, and plot it using using
the language of your choice.
(b) At what value of time t is the tension a maximum, and at what value is it a
minimum? Do these results make sense?
19. In this chapter, we saw that for small amplitude oscillations, the period of a pendulum
is τ = 2π L/g.
p
(a) Show that in the more general case of oscillations with any amplitude θ0 , the
period is given by the expression:
Z θ0
dθ
τ = 8L/g
p
√
0 cos θ − cos θ0
(b) By expanding the integrand in a series of the angle θ, show that the integration
can be performed term-by-term to obtain the following expression:
θ2
τ = 2π L/g 1 + 0
p
16
i i
i i
i i
derive
x(t) = Ae−γt cos (ωd t − φ)
q
for the case of underdamped motion, where ωd = ω02 − γ 2 and γ < ω0 . What is the
mathematical relationship between the constants A0 , A1 and the constants A, φ?
26. Use the weak damping approximation γ << ω0 to compute x(t) for the case of under-
damped motion. What is the resulting equation in the limit of γ → 0?
i i
i i
i i
27. Show that x(t) = teλt is a solution to the equation ẍ + 2γ ẋ + ω02 x = 0, in the case of
critically damped oscillations.
d2 x dx
+ 6 + 8x = 10 cos(2t)
dt2 dt
At time t = 0 the particle is at the origin and at rest.
(a) Find the transient and steady-state solutions for the position and velocity of the
oscillator as functions of time t. Do this by hand, and also by using Mathematica
or Python.
(b) Find the amplitude and period of the oscillator after a long time has elapsed.
(c) Find the rate of change of the kinetic energy dT/dt as a function of time using
Mathematica or Python and plot this quantity over many periods of oscillation.
Discuss the shape of the resulting graph at different intervals of time t.
30. The position x(t), y(t) of a particle on the xy-plane is given by x(t) = A cos(ω1 t + φ1 )
and y(t) = B cos(ω2 t + φ2 ). Prove that if the ratio ω1 /ω2 is a rational number, then
the particle moves in a closed curve.
31. The position x(t), y(t) of a particle on the xy-plane satisfies the equations
d2 x
= −8y
dt2
d2 y
= −2x
dt2
At time t = 0, the particle is at the position (1, 0) and at rest. Find the position and
velocity as a function of time t. Do this problem by hand and also using a computer.
32. A particle with mass m = 1 kg moves on the xy-plane where the potential energy
function is V (x, y) = x2 + y 2 + xy + 3. The particle is initially at the point (1, 1) and
moves with a speed of 2 m/s along the positive x-axis.
(a) Find the position as a function of time t.
(b) Find the period of the motion.
i i
i i
i i
Sections 6.7-6.8: Energy Resonance and the Quality Factor for Driven Oscil-
lations and Electrical Circuits
33. An electrical circuit consists of a 10 Ω resistor, a 80 mH inductance, a 10 µF capacitor,
and an AC voltage of the form V = 100 cos(10t) (where V is measured in volts and
time in seconds) connected in series.
(a) Find the current I(t) in the circuit at any time t. Plot this function and identify
the steady state part and the transient part of the graph.
(b) Plot the electrical energy stored in the capacitor as a function of time, and discuss
the shape of this graph.
(c) Find the quality factor Q in this circuit.
(d) Find the average power P = V I dissipated in the resistor R, between t = 0.2 s
and t = 0.4 s.
34. Give examples of mechanical and electrical systems that have a very low Q-factor and
a very high Q-factor. What is the physical meaning of these numbers for the oscillators
they describe?
Sections 6.9: Superposition Principle, Fourier Series
35. Use the symbolic capabilities of Mathematica or another CAS, to:
(a) Find analytical expressions for the Fourier coefficients for the periodic square
wave function shown in Figure 6.17.
(b) Plot the sum of the first 10 terms in the Fourier series and discuss how the plot
compares with the periodic square wave.
(c) Repeat part (b) for the first 100 terms of the Fourier series and discuss the
results.
3.0
2.5
2.0
f(x)
1.5
1.0
0.5
0.0
-5 0 5
Position x
Figure 6.17: Problem 6.35. Piecewise continuous function, f (x).
i i
i i
i i
CHAPTER 7
In the next chapter, we will develop a new way of finding a system’s equations of motion. The
new method will involve a branch of mathematics called the calculus of variations, which
we introduce in this chapter. The calculus of variations focuses on finding the extrema of a
mathematical object called a functional. The foundations of the calculus of variations were
laid by Isaac Newton and Gottfried Wilhelm Leibniz (1646–1716) and was further developed
by the brothers Jakob and Johann Bernoulli (1655–1705 and 1667–1748, respectively). The
first major contributions to the calculus of variations were made by Leonhard Euler (1707–
1783), Joseph-Louis Lagrange (1736–1813), and Pierre-Simon Laplace (1749–1827). In this
chapter, we will introduce the calculus of variations and use it to solve problems related to
physics.
209
i i
i i
i i
r2
y2
r1
y1 y=y(x)
x1 x2
Figure 7.1: A function y(x) which describes a path between two points, r1 = (x1 , y1 ) and
r2 = (x2 , y2 ).
i i
i i
i i
where r1 and r2 denote the location of the first and second point, respectively. As you may
recall from your introductory calculus course, it is common to factor out the dx2 term from
ds in order to obtain:
Z x2 q
L= 1 + y 0 (x)2 dx (7.2.2)
x1
where we used:
dy
dy = dx = y 0 (x)dx (7.2.3)
dx
In the introductory paragraphs of this chapter we said, “[the] calculus of variations provides
a means of finding the function which extremizes the quantity of interest.” In this case, the
“quantity of interest” is the length L of the curve y(x) joining the points r1 and r2 . Notice
that the length of the curve depends on the function y(x) (or more specifically, its derivative),
hence our “quantity of interest” is a function of the function y(x) and, therefore, we can
write the length of the curve as L = L(y). The problem specifically states that we want to
find the path which minimizes the distance between the points r1 and r2 . In other words,
we want to find the function y ∗ (x) such that L(y ∗ ) is a minimum value of L. As we will
see, the calculus of variations will provide us with a means of how to find y ∗ (x).
Before showing how to solve the above problem, it will be useful to establish some
notation. Notice that (7.2.2) takes the form:
In the case of (7.2.2), the integrand f [y(x), y 0 (x); x] = 1 + y 0 (x)2 is a function of only
p
y 0 (x). The use of square brackets denotes the functional nature of f , i.e., that, loosely
speaking, its arguments are functions. The semicolon is used to separate out the independent
variable in the argument of f from the functions. In other words, f depends on functions
y(x) and y 0 (x) but can also contain terms of those functions’ independent variable, x. It
should be noted that there are no hard rules for these two notations, square brackets and
semicolons, parentheses and a comma could have been used instead. Many of the problems
from the calculus of variations that we will study will involve integrands that are functions of
y(x), y 0 (x), and x. However, in this case, because both y and y 0 depend on x, the integrand
is actually only a function of one variable, x. Later on in this chapter, we will explore cases
where the integrand is a function of multiple variables. Furthermore, J is sometimes referred
to as a functional. A functional is a mathematical object that takes a function as an input
and returns a scalar as an output. In the case of (7.2.2), the functional L can be rewritten
as L(y), and takes the equation of the curve y(x) as an input and returns as the length of
the curve (a scalar) as the output.
Next, we will develop a method for finding the function y(x) that is an extremum of the
functional J(y). Or, continuing with our example, we will learn how to find the curve y ∗ (x)
that minimizes the path length between two points on a plane.
i i
i i
i i
stationary. It will be helpful to think of the “shortest distance between two points on a
plane” example from the previous section as we work through the following derivation.
However, keep in mind that J is not necessarily a path length. Let us return to Figure 7.1
for a visualization; however, this time we will start with the true stationary function y(x)
and vary it by adding a small perturbation η(x) such that η(x1 ) = η(x2 ) = 0. A perturbation
is a small additive deviation from a quantity. Hence, the perturbation results in a function
Y (x) = y(x) + aη(x), shown in Figure 7.2, which no longer makes the functional J stationary
because of the perturbation. The parameter a provides a measure of the amount that y(x)
is perturbed.
The functional J becomes:
Z x2
J = f [Y, Y 0 ; x]dx (7.3.3)
x
Z x1 2
= f [y + aη, y 0 + aη 0 ; x]dx (7.3.4)
x1
In other words, J = J(a) because the integration will “remove” J’s dependence on x.
In effect, we have reduced our problem to a function of a single variable a, for which we
know how to find its extrema by using introductory calculus. Because we know that the
i i
i i
i i
Y(x)=y(x)+aη(x)
y2
y1
y(x)
x1 x2
Figure 7.2: The function η(x) perturbs y(x), the extremum of (7.3.2), resulting in the
functional J not being stationary. The parameter a describes the size of the perturbation
from the extremum y(x).
extremum exists when a = 0, we need to show that ∂J/∂a equals zero when a = 0. As we
will see, the result of this work will be a condition for which y(x) makes J stationary.
Next we take the derivative of J with respect to a:
Z x2
∂J ∂
= f [Y, Y 0 ; x]dx (7.3.5)
∂a ∂a x1
Z x2
∂f ∂Y 0
∂f ∂Y
= + dx (7.3.6)
x1 ∂Y ∂a ∂Y 0 ∂a
∂Y ∂Y 0 dη
=η and =
∂a ∂a dx
Therefore,
Z x2
∂J ∂f ∂f dη
= η+ dx (7.3.7)
∂a x1 ∂Y ∂Y 0 dx
Next, we integrate the second term in (7.3.7) by parts using:
∂f dη
u= dv = dx
∂Y 0 dx
d ∂f
du = dx v =η
dx ∂Y 0
such that,
x2 Z
Z x2 xx
∂f dη ∂f d ∂f
dx = η(x) − η(x)dx (7.3.8)
0
∂Y dx ∂Y 0 dx ∂Y 0
x1 x1
x1
The first term on the right-hand side of (7.3.8) is zero because η (x1 ) = η (x2 ) = 0 as can be
seen in Figure 7.2. This leaves the second term remaining in (7.3.8), which will be inserted
into (7.3.7) to obtain
i i
i i
i i
Z x2
∂J ∂f d ∂f
= − η(x)dx (7.3.9)
∂a x1 ∂Y dx ∂Y 0
Next, we need to evaluate (7.3.9) at a = 0, resulting in
Z x2
∂J ∂f d ∂f
= − η(x)dx = 0 (7.3.10)
∂a x1 ∂y dx ∂y 0
where we have set the derivative equal
R to zero to ensure that J is stationary.
Equation (7.3.10) is of the form, h(x)η(x)dx = 0, and must be true for any continuous
function η(x). It can be shown that if h(x)η(x) = 0 for any η(x) then h(x) = 0 for all x.
R
which provides a means for finding the function y(x) that makes J stationary. Equation
(7.3.11) is also sometimes referred to as the first form of the Euler equation and provides
the promised condition for finding the extremum y because, as we will see, (7.3.11) will
produce a differential equation for y(x).
Next, we will use the Euler equation to find the path of the shortest distance between
two points on the Cartesian plane.
Example 7.1: The shortest path between two points on the Cartesian plane
Find the path that has the shortest length between two points (x1 , y1 ) and (x2 , y2 ) on
the Cartesian plane.
Solution:
We start with (7.2.2),
Z x2 q
L= 1 + y 0 (x)2 dx
x1
where we see that f = 1 + y 0 (x)2 . Next, we compute the derivatives needed for the Euler
p
equation:
∂f ∂f y0
=0 and = .
∂y 0
q
∂y
1 + (y 0 )2
Because the first term of the Euler equation is zero, we have
d y0 =0
q
dx 2
1 + (y )
0
or, in other words, the term within the parentheses must be constant in x. Setting that
constant equal to c produces
i i
i i
i i
y0
q =c
1 + (y 0 )2
which leads to:
c
y0 = √ ≡m
1 − c2
The term on the right-hand side of the above equation is also a constant which we
defined to be m. As promised, the Euler equation leads to a differential equation from
which y can be found. It is easy to check that the differential equation y 0 = m results in:
y(x) = mx + b
where b is the constant of integration. Hence, the path with the shortest distance between
two points in the Cartesian plane is a straight line. No surprise. Computer algebra systems
sometimes have algorithms for handling problems involving the calculus of variations. The
following Mathematica code creates the Euler equation for this example.
<< VariationalMethods`
y 00 [x]
OUTPUT: − == 0
(1+y0 [x]2 )
3/2
y 00 [x]
DSolve − == 0, y[x], x
(1+y0 [x]2 )
3/2
df ∂f dy ∂f dy 0
= + (7.4.1)
dx ∂y dx ∂y 0 dx
df ∂f ∂f
= y0 + y 00 0 (7.4.2)
dx ∂y ∂y
i i
i i
i i
∂f df ∂f
y 00 = − y0 (7.4.3)
∂y 0 dx ∂y
In addition, we can compute:
d 0 ∂f ∂f d ∂f
y = y 00 0 + y 0 (7.4.4)
dx ∂y 0 ∂y dx ∂y 0
Finally, we insert (7.4.3) into (7.4.4),
d ∂f df ∂f d ∂f
y0 0 = − y0 + y0 (7.4.5)
dx ∂y dx ∂y dx ∂y 0
d ∂f ∂f df
= y0 − + (7.4.6)
dx ∂y 0 ∂y dx
The term in the square brackets in (7.4.6) is Euler’s equation, and if y(x) makes J
stationary, then that term is zero according to the first form of the Euler equation. Therefore
we have:
d 0 ∂f
f −y =0 (7.4.7)
dx ∂y 0
Using the same kind of argument as we did in Example 7.1, we can set the term inside the
parentheses equal to a constant and therefore:
The second form of the Euler equation is often useful when the dependent variable, x
does not explicity appear in f , i.e., ∂f /∂x = 0. An example of such a case is when finding
the mathematical form of a geodesic on a sphere (see below).
i i
i i
i i
(x1,y1)
x
F
(x2,y2)
Figure 7.3: The solution of the brachistochrone problem will provide the formula for the
dotted path which minimizes the time of travel for a particle moving between the two points
(x1 , y1 ) and (x2 , y2 ) in a constant conservative force field, F = mg ĵ.
In our previous example,pwe used ds = 1 + y 0 (x)2 dx, however, we are going to find it more
p
i i
i i
i i
(x1,y1)
x
a
θ
(x2,y2)
y
Figure 7.4: The cycloid is the path that makes (7.5.3) stationary. The path is parameterized
by the variable θ, which is defined as shown in the figure above.
∂f x0 1
= =√ (7.5.6)
0
y (1 + x ) 2a
p
∂x 02
√
where a is a constant. The choice of using 1/ 2a as our constant may seem strange, but it
will simplify the mathematics later on in the calculation. Next, we will square (7.5.6) and
solve for x0 to obtain,
r
y
x0 = (7.5.7)
2a − y
which is a differential equation whose solution will give us the function x(y) that makes
(7.5.3) stationary. To solve for x(y), we need to solve the integral,
Z r
y
x= dy (7.5.8)
2a − y
which can be done using the substitution:
y = a (1 − cos θ) (7.5.9)
The result of the integration is that,
i i
i i
i i
Solution:
The code to create the plot of the cycloid in Mathematica is shown below.
a = 2;
7.5.2 Geodesics
A geodesic is the shortest possible line between two points on a curved surface. In this
example, we will find the geodesic on a sphere of radius R. We begin by noting that the
infinitesimal displacement along the surface of the sphere (dr = 0) is:
q p
ds = R dθ2 + sin2 θdφ2 = R θ02 + sin2 θ dφ
where we used θ0 = dθ/dφ. In this case, we are going to find the geodesic as a curve θ(φ).
If we use the Earth as our sphere, then finding θ(φ) is similar to describing the latitude of
the points along the geodesic as a function of their longitude.
The length integral we are attempting to make stationary is,
Z Z φ p
L= ds = R θ02 + sin2 θ dφ (7.5.11)
φ0
dφ a csc2 θ
=√
dθ 1 − a2 csc2 θ
which when integrated gives,
i i
i i
i i
cot θ
φ = sin−1 +α (7.5.15)
β
where β 2 = (1 − a2 )/a2 and α is a constant of integration. Finally, we can rewrite (7.5.15)
as:
x =R sin θ cos φ
y =R sin θ sin φ (7.5.18)
z =R cos θ
z = Ay − Bx (7.5.19)
which is the equation of a plane passing through the center of the sphere. Because we are
restricted to the surface of a sphere of radius R, the geodesic is the path made by the
intersection of the plane with the sphere’s surface. Hence, the geodesic is the great circle
that lies at the intersection of the plane and the surface of the sphere. The great circle
between two points is the path often used for air and sea navigation, because it provides
the path with the shortest distance between those points.
i i
i i
i i
z
(y2,z2)
ds
(y1,z1)
x
Figure 7.5: The surface of revolution formed by a line joining the points (y1 , z1 ) and (y2 , z2 ).
this case, r would be the distance y from the z-axis to the curve defining the surface of
revolution. The height of the cylinder ds is the differential distance along the curve z(y),
ds = dy 2 + dz 2 . Therefore, our functional for the area is:
p
Z y2 p
A= 2πy 1 + z 02 dy (7.5.20)
y1
We want to find z(y) that minimizes the area of the surface of revolution. Therefore
√ we
factored dy out of ds to produce (7.5.20), where z 0 = dz/dy. Next, we insert f = y 1 + z 02
into the first form of the Euler equation (7.3.11),
yz 0
d
√ =0 (7.5.21)
dy 1 + z 02
Note that ∂f /∂z = 0. Therefore we can write,
yz 0
√ =a (7.5.22)
1 + z 02
where a is a constant. Solving (7.5.22) for z 0 we obtain,
dz a
=p (7.5.23)
dy y − a2
2
i i
i i
i i
1
a cosh =3 (7.5.26)
a
which must be done using numerical approximations. The following Mathematica code
numerically solves (7.5.26) for a.
i i
i i
i i
The default for RevolutionPlot3D is to plot a function fz (t) which describes the height
of the curve above the xy-plane, where t is the radius from the z-axis. The command
then rotates the curve of fz about the z-axis. However, our function describes the height
of the curve above the xz-plane. In order to get the shape we expect, we need to tell
Mathematica to revolve about what it believes to be the x-axis by using the option,
RevolutionAxis→ {1, 0, 0}. This is the type of technical issue one occasionally encounters
when working with software. If you had plotted (7.5.25) using the defaults, you would have
obtained a plot that was clearly incorrect. The plot would have looked like a bowl and the
axes would have had the wrong ranges as compared to what was described in the problem
(notice our solution has rings of radius 3, one at z = 1 and the other at z = −1). This is
another reminder that one can’t simply use computer algorithms as black boxes. You need
to make sure that you understand the algorithm and check to see that the output makes
sense.
where there are n dependent variables yi . In this case, it is often convenient to write f =
f [yi , yi0 ; x], where i = 1, 2, . . . , n.
We want to find a set of functions {yi (x)} (i = 1, 2, . . . , n) that makes the functional
Z
J = f yi (x), yi0 (x); x dx (7.6.2)
stationary. To find {yi (x)}, we follow the work done in Section 7.2. We write Yi (α, x) =
yi (x) + αηi (x), insert Yi into (7.6.2), differentiate the resulting J with respect to α, and
then set that derivative equal to zero. The result is,
i i
i i
i i
Z n
x2 X
∂J ∂f d ∂f
= − ηi (x)dx (7.6.3)
∂α x1 i=1 ∂Yi dx ∂Yi0
Each perturbation ηi (x) is independent of one another. Therefore, in order for the derivative
in (7.6.3) to be zero at a = 0, we must have,
In other words, we will have an Euler equation for each dependent function yi . As before,
each Euler equation will produce an ordinary differential equation. However, these ODEs
are often coupled and can be difficult to solve in closed form.
For an example of f depending on multiple dependent variables, we can reconsider the
first problem we studied, the shortest path between two points on the Cartesian plane. When
we first solved that problem, we assumed y = y(x). However, if we wanted to consider all
possible paths, then we would need to consider curves which are parameterized using a
variable t,
where x0 = dx/dt and y 0 = dy/dt. Using f = x02 + y 02 , and the Euler equations:
p
d ∂f ∂f d ∂f ∂f
− =0 and − =0
dt ∂x0 ∂x dt ∂y 0 ∂y
we find,
x0 y0
=a and =c (7.6.6)
x02 + y 02 x02 + y 02
p p
where a and c are constants. Note that we used the fact that ∂f /∂x = ∂f /∂y = 0, in order
to set ∂f /∂x0 and ∂f /∂y 0 equal to the constants a and c, respectively. By dividing the two
equations in (7.6.6), we get,
y0 dy c
0
= = =m (7.6.7)
x dx a
Integrating (7.6.7), we find y = mx + b, where b is the constant of integration. In other
words, we again get the equation for a line.
i i
i i
i i
Common problems in the calculus of variations involve finding the path of the shortest
distance between two points, finding the path that minimizes the time of travel between
two points, and finding the minimum surface area for a surface of revolution, to name a
few.
The function y(x) can be found using Euler’s equation:
d ∂f ∂f
0
− =0
dx ∂y ∂y
Euler’s equation will produce a differential equation, which can be solved for y(x). An
alternative form of the Euler’s equation is,
∂f
f − y0 = constant
∂y 0
which can be used when ∂f /∂x = 0.
Problems involving finding the extrema of a functional which depends on several depen-
dent variables {yi (x)}, for i = 1, 2, . . . , n, of the form,
Z x2
J= f [yi (x), yi0 (x); x]dx
x1
i i
i i
i i
(x1,y1) (x2,y2)
θ1 θ2
(xr,yr)
x
Figure 7.6: Problem 7.5: The path of a light ray in the xy-plane reflecting off of a mirror
in the xz-plane.
stationary with y 0 = dy/dx. The path must connect the points (0, 0) and (1, 1).
4. Find the shortest path between two points on the surface of a cylinder with radius R.
Consult Problem 2, for how to set up the necessary functional.
5. Consider the light ray (solid line) reflecting off of the surface of a plane mirror at the
point (xr , yr ) as shown in Figure 7.6. The surface of the mirror is in the xz-plane,
and the light ray is in the xy-plane. Using Fermat’s principle from Problem 1, prove
the law of reflection. In other words, show that the angle of incidence θ1 is equal to
the angle of reflection θ2 in the figure below. Although not necessary, you may find
it helpful to set x1 = 0 and yr = 0 as suggested in Figure 7.6. To solve the problem,
compute the time it takes to travel along the path in Figure 7.6, then show that the
time is a minimum when θ1 = θ2 . You may start by assuming that the point (xr , yr )
is in the same vertical plane as (x1 , y1 ) and (x2 , y2 ). If you were not to make that
assumption, then explain how you would need to change your solution in order to
prove the law of reflection.
i i
i i
i i
(x1,y1) y
θ1
(xr,yr) n1
n2
x
θ2
(x2,y2)
Figure 7.7: Problem 7.6: The path of a light ray in the xy-plane leaving a medium with an
index of refraction n1 and entering a medium with an index of refraction n2 .
6. Consider the light ray shown in Figure 7.7. The light ray starts in a medium with a
spatially uniform index of refraction n1 and, at the point (xr , yr ) enters a new medium
with a spatially uniform index of refraction n2 . Using Fermat’s principle from Problem
1 to minimize the time of travel between the two points (x1 , y1 ) and (x2 , y2 ), prove
Snell’s law, n1 sin θ1 = n2 sin θ2 and that the point of refraction (xr , yr ) is coplanar
with the points (x1 , y1 ) and (x2 , y2 ). As before, the index of refraction is defined as
the ratio n = vc , where c is the speed of light in vacuum and v is the speed of light
inside the medium.
where y 0 = dy/dx.
Hint: Rewrite the functional in polar coordinates and use Mathematica’s Variational-
Methods package to get Euler’s equations.
9. Find the geodesics on the surface of the cone whose equation in cylindrical coordinates
is z = cr.
10. Imagine a ball rolling down the first half of the first loop of a cycloid (from x = 0 to
x = πa), defined by the equations:
y = a (1 − cos θ) x = a (θ − sin θ)
Show that the time of travel from x = 0 to x = a is equal to π a/g. Show that this
p
i i
i i
i i
11. For a little bit of foreshadowing of upcoming chapters, consider the functional,
1 1
Z t
J(x) = mx02 − kx2 dt0
t0 2 2
where m and k are positive constants, x0 = dx/dt0 and t0 is the independent variable
(and not a derivative). Find the function x(t) that makes the functional stationary.
Does this look familiar?
13. Repeat Problem 4, however, this time don’t assume θ = θ(z). Instead, you should
parameterize the curve such that z = z(t) and θ = θ(t), where t is the parameter of
the curve.
14. Find the shortest path on the conical surface z = z0 − c x2 + y 2 , where z0 and c are
p
constants.
15. Find the extrema of the functional,
Z π
J(y, z) = y 02 + z 0 + yz dx
0
where y 0 = dy/dx and z 0 = dz/dx, subject to the conditions, y(0) = 0, y(π) = 1, z(0) = 0,
and z(π) = 1.
i i
i i
i i
CHAPTER 8
In this chapter, we will present an alternative way of formulating physics problems. This
new method was first published by French-Italian mathematician Joseph-Louis Lagrange
(1736–1813) in 1788 and is therefore called Lagrangian dynamics. Lagrange’s formulation
uses a function called the Lagrangian, which is found using the particle’s energy. As we
will see, a particle will move in such a way as to make the integral of the Lagrangian
with respect to time stationary. The integral, called an action integral, is a functional (a
mathematical formula that maps a function to a number). Hence, the particle’s path can
be described using a function that is an extremum of a functional (the action integral). In
the last chapter, we learned how to find the functions that are the extrema of a functional
by using the Euler equation. In this chapter we will learn how the Euler equation can be
used to find a system’s equation of motion.
One advantage of the Lagrangian formulation is that it is the same, regardless of the
coordinate system being used. Lagrangian dynamics uses the Euler equation to derive a
particle’s equation of motion and the Euler equation takes the same form regardless of
the coordinate system being used. As we saw in Chapter 3, Newton’s second law takes
different forms depending on the coordinate system, because the formula for the acceleration
depends on the coordinate system. Another advantage of Lagrangian dynamics is that it
eliminates the need of knowing constraint forces. A bead constrained to move along a wire
is a common example involving constraint forces. The force keeping the bead on the wire
is often unknown, or very difficult to find. Lagrangian mechanics doesn’t need to know the
formula of the constraint force in order find the particle’s motion.
The final “advantage” we’d like to mention is that Lagrangian mechanics is simply
elegant! As we will see, visual inspection of the Lagrangian can reveal what quantities are
conserved in the system. We will also learn how to use the Lagrangian to derive another
quantity called the Hamiltonian, which can also be used to derive a particle’s equations of
motion. Furthermore, the Hamiltonian will provide us with a new way of thinking about the
phase space, where we can use pictures to describe the possible behaviors of a system. The
Hamiltonian’s value extends well beyond the field of classical mechanics and is the central
quantity used to describe a particle’s wave function in quantum mechanics.
229
i i
i i
i i
mẍ + kx = 0 (8.1.3)
You may recall that normally we define ω02 = k/m for the SHO, but for now we will leave
the equation of motion as it appears above.
There is another method of finding (8.1.3) that doesn’t directly involve Newton’s second
law. Notice that,
d ∂T d
= (mẋ) = ṗ = mẍ (8.1.4)
dt ∂ ẋ dt
Therefore, we can obtain (8.1.3) by writing,
d ∂T ∂V
+ =0 (8.1.5)
dt ∂ ẋ ∂x
Equation (8.1.5) is similar to the Euler equation in appearance, except that the functions
appearing in each derivative are different. We can make (8.1.5) look like the Euler equation
by defining a new function, the Lagrangian,
L ≡ T −V (8.1.6)
Notice that V does not depend on ẋ , and T does not depend on x, therefore ∂T /∂x = 0
and ∂V /∂ ẋ = 0 and we can write:
d ∂T d ∂L
= (8.1.7)
dt ∂ ẋ dt ∂ ẋ
and
∂V ∂L
=− (8.1.8)
∂x ∂x
And therefore we obtain from (8.1.5) the so-called Euler-Lagrange equation,
i i
i i
i i
along the pathR of the oscillator, x(t). In other words, the path of the oscillator is the one
that makes L(x, ẋ; t)dt stationary. We will formalize this statement in the next section.
However, the important point here is that we did not use Newton’s second law to derive
the equation of motion for the SHO. Instead of asking what forces are acting on the system
(as is done using Newton’s second law), we asked what behavior makes the integral of the
Lagrangian stationary. The Lagrangian depends only on the kinetic and potential energies
of the system. Note, however, that the Lagrangian is not the total mechanical energy of
the system. The total mechanical energy is the sum of the kinetic and potential energies,
while the Lagrangian is the difference between them.
Although we have demonstrated the Euler-Lagrange equation for the SHO, as we will
see in later sections, the Euler-Lagrange equation can be used to find the equation of motion
for any system. Before generalizing the work done in this section to other systems, we need
to discuss coordinate systems.
i i
i i
i i
y
θ ℓ
x
m
Figure 8.1: A simple plane pendulum composed of a mass m attached to a massless rigid
rod of length ` which makes an angle θ with the vertical (dashed line).
Next, suppose that there are constraints on the particle’s motion. Each constraint will
relate one quantity to another. For example, each particle may be constrained to move on
the xy-plane. Regardless of their specific nature, constraints can relate two or more of the
quantities to each other or they can impose a value for on particular coordinage (e.g. z = 0).
In other words, not all 3N quantities are necessarily independent. If there are m equations
of constraint, then there are s = 3N − m degrees of freedom for the system. The number
of degrees of freedom is the number of independent quantities that are needed to specify
the state of the system. In other words, we will need s independent generalized coordinates
in order to describe the state of the system. As we will see, the Euler-Lagrange equations
will produce s second order differential equations which, when solved, will give qj (t) for
j = 1, 2, . . . , s.
The following example will help clarify constraints and degrees of freedom. Consider the
simple plane pendulum shown in Figure 8.1.
In the case of the simple plane pendulum there is one particle, the mass m, and there-
fore N = 1. Generally speaking, the position of the mass would require three coordinates,
x, y, and z. However, because we are working with a plane pendulum, we have our first
equation of constraint, z = 0. Next, we know that the mass is constrained to move along
a circle of radius `. Therefore, the next equation of constraint is x2 + y 2 = `2 . We see that
we have two equations of constraint, therefore we have s = 3 (1) − 2 = 1 degree of freedom.
To describe the system, we’d like to work with a generalized coordinate system. The most
obvious one to choose is polar coordinates. In this case, we’ll choose q = θ, where θ is the
angle between the pendulum and the dashed line in Figure 8.1. We can write the Cartesian
coordinates in terms of the generalized coordinate θ using:
i i
i i
i i
1 1
T = m ẋ2 + ẏ 2 = m`2 q̇ 2 (8.2.3)
2 2
d ∂L ∂L
− =0 (8.2.6)
dt ∂ q̇ ∂q
will produce the equation of motion. Inserting (8.2.5) into (8.2.6) gives the simple plane
pendulum equation,
g
q̈ + sin q = 0 (8.2.7)
`
The path taken by a particle from the point r1 to the point rR2 during the time interval
t
t1 to t2 is the one that makes the action integral S = t12 Ldt stationary.
Hamilton’s principle is a cornerstone not just of classical physics, but also of quantum
physics. From Chapter 7, we know that the condition for S to be stationary is
i i
i i
i i
Q(t)=q(t)+aη(t)
q2
q1
q(t)
t1 t2
Figure 8.2: The function q(t) is the path that makes Ldt stationary and Q(t) is a pertur-
R
Euler equation in Chapter 7, and therefore we want to find a function q(t) that makes the
action integral,
Z t2
S= L (q, q̇) dt (8.3.2)
t1
stationary. Following the argument from Chapter 7, we will perturb q(t) using a function
η(t), which is zero at the limits of integration, as shown in Figure 8.2. The function Q(t) is
the perturbation of q(t) and the perturbation has a size a, which is constant.
Next, we will differentiate S with respect to a and the extrema function q(t) should be
found when,
∂S
=0 (8.3.3)
∂a
a=0
From Chapter 7, we know that the above derivative is zero when,
d ∂L ∂L
= (8.3.4)
dt ∂ q̇ ∂q
If we substitute,
1
L = mq̇ 2 − V (q) (8.3.5)
2
into (8.3.4), then we have:
d ∂V
(mq̇) = − (8.3.6)
dt ∂q
The right-hand side of (8.3.6) is the q component of the force acting on the particle.
Likewise, the left hand side of (8.3.6) is the time derivative of the momentum. Therefore,
we see that the path q(t) that makes the action integral S stationary is the one that
also solves Newton’s second law. Therefore, Hamilton’s principle does provide a means of
finding a system’s equation of motion. Note that the above work is not a proof of Hamilton’s
principle, but rather a demonstration of how it is capable of producing a system’s equation
of motion.
i i
i i
i i
Before we move onto some examples of using the Lagrangian to find the equation of
motion for a system, it is important that we point out some additional terminology that
will become useful later. First, notice V = V (qj ) and T = T (q̇i ). Hence, in our example
∂L ∂L ∂V
= mq˙j and =− (8.3.7)
∂ q̇j ∂qj ∂qj
Term ∂L/∂ q̇ is sometimes called the j th component of the generalized momentum although
it is not always a linear momentum (mass times velocity). Likewise, ∂L/∂qj is sometimes
called the j th component of the generalized force; again it is not always a force, but it
acts like one. With these terms in mind, (8.3.4) can be reworded as: “the generalized force
is equal to the time rate of change of the generalized momentum.” It is a “generalized”
Newton’s second law!
stationary.
Because there are three degrees of freedom, i.e., three functions that need to be found
in order to make S stationary, there will be three Euler-Lagrange equations. Before finding
the Euler-Lagrange equations for this problem, first note that:
i i
i i
i i
d ∂L d
= (mẋ) = mẍ
dt ∂ ẋ dt
d ∂L d
= (mẏ) = mÿ (8.4.3)
dt ∂ ẏ dt
d ∂L d
= (mż) = mz̈
dt ∂ ż dt
By equating the terms in parentheses in each line above, we see that ∂L/∂ ẋ is the
x-component of the particle’s momentum (similar for the ẏ and ż derivatives). In this
case, the generalized momentum is the same as the linear momentum you learned about
in your introductory physics class. Likewise,
∂L ∂V
=− = Fx
∂x ∂x
∂L ∂V
=− = Fy (8.4.4)
∂y ∂y
∂L ∂V
=− = Fz
∂z ∂z
We can see that ∂L/∂x is the x-component of the force acting on the particle (similar
for the y and z derivatives). In this case, the generalized force is equal to the force acting
on the particles. Inserting the derivatives into the Euler-Lagrange equation gives:
d ∂L ∂L
− =0 ⇒ mẍ = Fx
dt ∂ ẋ ∂x
d ∂L ∂L
− =0 ⇒ mÿ = Fy
dt ∂ ẏ ∂y
d ∂L ∂L
− =0 ⇒ mz̈ = Fz
dt ∂ ż ∂z
The result of the Euler-Lagrange equations is simply Newton’s second law.
While Example 8.1 may seem elementary, it points out a few important things. First,
the number of Euler-Lagrange equations will equal the number of degrees of freedom in
the system. Second, Example 8.1 demonstrates the connections between the generalized
momentum and generalized force to the physical quantities momentum and force with which
you are already familiar. Of course, to progress further in this problem and find x(t), y(t),
and z(t), you would need to know V (x, y, z).
i i
i i
i i
R
x1
x2
m1
m2
We use the horizontal dashed lined passing through the center of the Atwood machine
as the origin for our coordinate system. The coordinates x1 and x2 measure the distance
from the center of the Atwood machine to masses m1 and m2 , respectively. Because the
string is inextensible, the two coordinates x1 and x2 are not independent. For example,
as m1 falls down, m2 rises and vice versa. Therefore, there is one equation of constraint
among the coordinates:
` = x1 + x2 + πR (8.4.5)
The last term in the constraint is the length of the string that wraps around the pulley.
Of course an additional constraint is that the masses are each constrained to move only up
or down along a line, therefore y1 = y2 = z1 = z2 = 0, which are actually four constraints
(one on each coordinate) for a total of five constraints on the system (m = 5). Therefore
with two particles in the system N = 2, we have s = 3(2) − 5 = 1 degree of freedom for the
problem. We will rewrite the coordinates x1 and x2 in terms of the generalized coordinate
x using:
x1 = x (8.4.6)
x2 = ` − x − πR (8.4.7)
Notice that the velocities have the relationship, ẋ2 = −ẋ1 , and the accelerations have
the relationship, ẍ2 = −ẍ1 , which is what we would expect in this situation.
Next, we’ll find the kinetic and potential energies of the system. Notice that the pulley
has mass but the string does not. Therefore, the total kinetic energy includes the kinetic
energy of m1 , m2 , and the rotational kinetic energy of the pulley. You may recall from
your introductory physics course that rotational kinetic energy takes the form Trot = 21 Iω 2 ,
where I = 21 M R2 is the rotational inertia for the pulley, and ω = ẋ/R is the angular velocity
of the pulley. Therefore, the total kinetic energy for the system is:
1 1 1
T = m1 ẋ21 + m2 ẋ22 + Iω 2 (8.4.8)
2 2 2
2
1 1 1 ẋ
= m1 ẋ2 + m2 ẋ2 + I (8.4.9)
2 2 2 R
i i
i i
i i
The potential energy takes the form of the gravitational potential energy. Measuring
upward as the positive direction, we have:
1 1 1 I
L= m1 + m2 + ẋ2 + m1 gx + m2 g (` − x − πR) (8.4.12)
2 2 2 R2
Because s = 1, there is only one Euler-Lagrange equation for this system:
d ∂L ∂L
= (8.4.13)
dt∂ ẋ ∂x
d I
m1 + m2 + 2 ẋ = (m1 − m2 )g (8.4.14)
dt R
I
m1 + m2 + 2 ẍ = (m1 − m2 )g (8.4.15)
R
(m1 − m2 )g
ẍ = (8.4.16)
m1 + m2 + M
2
where we used I = 12 M R2 . From the generalized acceleration, we can get the acceleration
of each mass because according to (8.4.5), ẍ1 = ẍ and ẍ2 = −ẍ.
Example 8.2 demonstrates several ideas. First, drawing a picture of the system can help
identify constraints on the coordinates. Second, Example 8.2 reminds us to carefully identify
all kinetic and potential energies in the system. It would be easy to overlook the rotational
kinetic energy of the pulley if you thought the pulley is massless, a common assumption
in introductory physics. Finally, Example 8.2 shows us that although it might be useful to
describe the motion in generalized coordinates (in this case x), the solution of the problem
needs to address what is being asked by the problem. In this case, we needed to make sure
that we showed how to express the acceleration of each mass in terms of the generalized
coordinate.
Example 8.3: Pendulum supported by a rotating disk
A pendulum is suspended from a massless disk with a radius R that is rotating with
a constant angular velocity ω as shown in the figure below. The pendulum consists of a
mass m attached to one end of a rigid massless rod of length `. The other end of the rod
is attached to the edge of the rotating disk. Find the equation of motion of the mass m.
i i
i i
i i
y
ω
R x
ℓ
θ
m
Solution:
The coordinate system is defined in the figure where the coordinate axes are the dashed
lines. The mass is constrained to move in the xy-plane, hence z = 0 is our first equation
of constraint. There is an additional constraint which may not be immediately obvious.
However, if we write the location of the mass in terms of polar coordinates, then it will
become clear that this problem has only one degree of freedom. The coordinates of the
mass can be written as:
ẋ = Rω sin (ωt) + `θ̇ cos θ and ẏ = Rω cos (ωt) + `θ̇ sin θ (8.4.18)
Next, we will find the kinetic and potential energies. The kinetic energy T = 1
2 mv
2 can
be written using the velocity of the mass in the Cartesian plane, v = ẋî + ẏ ĵ,
1 1 1
T = mv 2 = mv · v = m ẋ2 + ẏ 2
2 2 2
The pendulum is in the Earth’s gravitational field, therefore the potential energy is
V = mgy. Inserting (8.4.17) and (8.4.18) into our equations for the kinetic and potential
energies results in the Lagrangian (after some algebra, see Problem 9),
1
L = m R2 ω 2 + `2 θ̇2 + 2R`ω θ̇ sin (θ − ωt) − mg (R sin (ωt) − ` cos θ) (8.4.19)
2
Inserting (8.4.19) into
d ∂L ∂L
− =0
dt ∂ θ̇ ∂θ
produces the equation of motion:
Rω 2 g
θ̈ + cos (θ − ωt) + sin θ = 0 (8.4.20)
` `
i i
i i
i i
Notice that (8.4.20) is the simple pendulum equation with an additional cosine term.
The cosine term is the drive term associated with the rotating disk.
Example 8.3 shows us another example of a problem which appears to have two degrees
of freedom, but can actually be described using only one. Another important point demon-
strated in Example 8.3 is that it is necessary to be careful with kinetic energies. Students
sometimes focus on potential energies, but this is a case where the kinetic energy is a bit
more difficult. It was important to include all of the components of the velocity in order to
correctly obtain the kinetic energy. Finally, if you solve Problem 9, you’ll see the value in
using trigonometric identities to cast the final solution in terms of something that is easier
to interpret.
m
R
y
r
θ
x
Solution:
We want to know something about the motion of the mass m. The system consists
of one particle N = 1 and in order to find the number of degrees of freedom s, we will
need to find the number of constraints. Cylindrical coordinates are a natural choice for
this problem. In those coordinates, we have two constraints. The first constraint is that
z = ar2 because the bead is constrained to move along the wire. The second constraint is
due to the rotation of the wire. We know that the wire is rotating about the z-axis at a
constant angular velocity, therefore, θ = ωt. Hence, we are left with one degree of freedom,
r.
Next we find the kinetic and potential energies. Using the velocity in cylindrical coor-
dinates, v = ṙr̂ + rθ̇θ̂ + żẑ, we find
1 1
T = mv · v = m ṙ2 + r2 θ̇2 + ż 2 (8.4.21)
2 2
The potential energy is found using V = mgz. By applying the constraints, we find:
i i
i i
i i
L =T − V
1
= m ṙ2 + r2 θ̇2 + ż 2 − mgz
2
1
= m ṙ2 + r2 ω 2 + 4r2 a2 ṙ2 − mgar2
2
Note that L = L (r, ṙ). To get the equation of motion for the bead, we can insert the
Lagrangian into the Euler-Lagrange equation,
d ∂L d m
= 2ṙ + 8r a ṙ = 1 + 4a r mr̈ + 8ma rṙ
2 2 2 2 2 2
dt ∂ ṙ dt 2
(8.4.23)
∂L
=mr rω + 4a ṙ − 2ga
2 2
∂r
Equating the two terms above gives:
We are interested in the case when the bead is traveling in a circle of radius R. There-
fore, inserting r̈ = ṙ = 0 and r = R into (8.4.24), we obtain
ω2
a=
2g
Example 8.4 demonstrates several important points. First, it demonstrates the need to
use the velocity in cylindrical coordinates. The velocity of a particle in polar, cylindrical,
and spherical coordinates was discussed in Chapter 3. It is important when computing the
kinetic energy that the velocity vector be expressed in the same coordinate system that you
are using to solve the problem. Second, although this problem didn’t specifically ask for the
equation of motion, we found it easier to answer the problem by finding it. Once we had the
general equation of motion for the bead (8.4.24), we inserted the specific conditions for the
motion of interest. Finally, Example 8.4 shows how the Lagrange formulation doesn’t need
to know all of the forces acting on the system. We do not know the mathematical form of the
force that keeps the bead on the wire. The aforementioned ignorance is no problem, all we
needed for the Lagrange formulation was how the value of the z coordinate was constrained.
i i
i i
i i
x2
x1 m
M θ
The figure above shows that we have created two coordinates, x1 and x2 . The coordi-
nate x1 is measured from the vertical axis and describes the location of the inclined plane
M . The coordinate x2 measures the position of the particle m from the top of the inclined
plane. Note that the two coordinates are not orthogonal to each other.
To calculate the Lagrangian, we need to find the kinetic energy of each object,
1 1
T = M v12 + mv22
2 2
where v12 = ẋ1 , but v22 6= ẋ22 . The particle’s velocity v2 relative to the vertical axis will
depend on the velocity of the inclined plane and therefore, v2 = ẋ1 + ẋ2 , where ẋ1 and ẋ2
are vectors that point along the directions of x1 and x2 in the figure, respectively. We can
find v22 = v2 · v2 by noting that the angle between ẋ1 and ẋ2 is θ,
1 1
T = M ẋ21 + m ẋ21 + ẋ22 + 2ẋ1 ẋ2 cos θ (8.4.28)
2 2
V = mg (` − x2 ) sin θ (8.4.29)
1 1
L = M ẋ21 + m ẋ21 + ẋ22 + 2ẋ1 ẋ2 cos θ − mg (` − x2 ) sin θ (8.4.30)
2 2
Note that we used ` as the length of the inclined plane and the height of the particle
above the horizontal axis is (` − x2 ) sin θ. Further note that the Lagrangian is of the form
L = L (x1 , x2 , ẋ1 , ẋ2 ). There are two degrees of freedom for this problem, therefore we will
have two Euler-Lagrange equations:
d ∂L ∂L
− =0
dt ∂ ẋ1 ∂x1
(8.4.31)
d ∂L ∂L
=0
−
dt ∂ ẋ2 ∂x2
Next we will compute the necessary derivatives,
i i
i i
i i
d ∂L ∂L
=M ẍ1 + m (ẍ1 + 2ẍ2 cos θ) =0
dt ∂ ẋ1 ∂x1
d ∂L ∂L
=m (ẍ2 + 2ẍ1 cos θ) =mg sin θ
dt ∂ ẋ2 ∂x2
i i
i i
i i
k1 =dt f (ti , xi )
∆t
k1
k2 =dt f ti + , xi +
2 2
∆t
k2
k3 =dt f ti + , xi + (8.5.2)
2 2
k4 =dt f (ti + ∆t, xi + k3 )
k1 k2 k3 k4
xi+1 =xi + + + + + O ∆t
5
6 3 3 6
The RK4 method computes the solution’s derivative once at each end point of an interval
and twice at the mid-point of the interval as shown in Figure 8.3. In Figure 8.3, the interval
is ∆t = ti+1 −ti and the circles represent the points in time where the derivative is evaluated.
The solid lines through the circles show the derivative of the solution (dashed line) at those
points.
i i
i i
i i
xi+1
xi
ti ti+Δt/2 ti+1
Figure 8.3: The derivatives (solid sloped lines) at the start, middle, and end points for a
given interval. Note that in this figure ∆t = ti+1 − ti .
To implement the RK4 algorithm, follow these steps in your programming language of
choice:
1. Define an empty array (or one with the initial condition as the first element) x, which
will be the solution x(t) of the ODE.
2. Define a time step ∆t which is small enough to produce the correct result.
3. Loop over an index i. During each iteration of the loop, the following values are
calculated: k1 , k2 , k3 , k4 , and xi . These are calculated in the order listed in (8.5.2)
from the top down.
As we will show in Example 8.6, second-order differential equations will require a modifica-
tion to the above process. For now, we focus on first-order ODEs.
Next, we will look at a sample RK4 algorithm, which we can use to solve a first-order
ODE. Consider the ODE,
ẋ = sin x (8.5.3)
with x(t = 0) = 1.0 m. Suppose that x measures the position of a particle whose velocity is
as described in (8.5.3). We want to solve for x(t) from t = 0 to t = 10 using a step size of
0.1. Algorithm 6 shows how to implement the RK4 method to solve (8.5.3) with the given
initial condition.
In Algorithm 6, we begin by defining the right-hand side of (8.5.3) as a function just like
we did with the Euler method. Next, we define the values of tmax = 10, dt = ∆t = 0.1, and
x0 = x(0) = 1.0. We then compute the number of steps n needed to integrate the solution
across the time range, and we set up arrays for t and x. The for-loop steps through the time
i i
i i
i i
range, computing the value of x at each time step idt. Notice that each of the k’s must be
updated for each step in the loop. The result is shown in Figure 8.4.
import numpy a s np
#d e f i n e f u n c t i o n
d e f f ( pos ) :
r e t u r n np . s i n ( pos )
#d e f i n e i n i t i a l c o n d i t i o n s and a r r a y s
x0 = 1 . 0 #i n i t i a l c o n d i t i o n
tmax = 1 0 . 0 #maximum time
dt = 0 . 1 #time s t e p
n = i n t ( tmax/ dt ) #number o f s t e p s
t = np . l i n s p a c e ( 0 , tmax , n ) #a r r a y t h a t s t o r e s t h e t i m e s
x = np . z e r o s ( n ) #a r r a y t h a t s t o r e s t h e s o l u t i o n
x [ 0 ] = x0
for i in range (1 , n ) : #l o o p f o r f i n d i n g t h e s o l u t i o n
k1 = f ( x [ i −1])
k2 = f ( x [ i −1]+k1 / 2 . 0 )
k3 = f ( x [ i −1]+k2 / 2 . 0 )
k4 = f ( x [ i −1]+k3 )
x[ i ] = x [ i −1] + dt ∗ ( k1 / 6 . 0 + k2 / 3 . 0 + k3 / 3 . 0 + k4 / 6 . 0 )
3.0
2.5
x, m
2.0
1.5
1.0
0 2 4 6 8 10
t, s
Figure 8.4: The solution to (8.5.3) as found by Algorithm 6.
As a demonstration on how to use the RK4 method to solve a second-order ODE, the
next example will numerically solve the equation of motion found in Example 8.3.
i i
i i
i i
Rω 2 g
θ̈ + cos (θ − ωt) + sin θ = 0
` `
Find θ(t) for t = 0 to t = 20 seconds using the values R = 0.2 m, ω = 3π/2 rad/s, ` = 1.0
m, g = 9.8 m/s2 , and initial conditions θ(0) = 1 rad and θ̇(0) = 0 rad/s.
Solution:
The first step in solving the problem is to create two first-order ODEs from the equation
of motion. Defining a new variable v = θ̇, we have as our new equations:
θ̇ =v
Rω 2 g (8.5.4)
v̇ = − cos (θ − ωt) − sin θ
` `
Next, we need to modify (8.5.2) for a system of two equations. For simplicity of nota-
tion, we will define the following functions:
Rω 2 g
f (t, θ, v) = v and cos (θ − ωt) − sin θ
h(t, θ, v) = −
` `
Then for a system of two equations, (8.5.2) becomes:
dt k1 j1
k1 =dt f ti + , θi + , vi + j1 =dt h (ti , θi , vi )
2 2 2
dt k1 j1 dt k1 j1
k2 =dt f ti + , θi + , vi + j2 =dt h ti + , θi + , vi +
2 2 2 2 2 2
dt k2 j2 dt k2 j2
k3 =dt f ti + , θi + , vi + j3 =dt h ti + , θi + , vi +
2 2 2 2 2 2
k4 =dt f (ti + dt, θi + k3 , vi + j3 ) j4 =dt h (ti + dt, θi + k3 , vi + j3 )
1 1
θi+1 =θi + (k1 + 2k2 + 2k3 + k4 ) vi+1 =vi + (j1 + 2j2 + 2j3 + j4 )
6 6
Algorithm 7 uses the RK4 method to solve (8.5.4) in Python.
The red (in the e-book) curled arrows in the code above denote a line continuation i.e.,
the text following a curly arrow is a continuation of the line above it. Notice that we used
f and g to denote our derivatives. Also, notice that the method is similar to what was done
for (8.5.3) except now we need to include the appropriate steps for each variable θ(ki ) and
v(ji ). A plot of θ(t) for the parameters listed in the problem appears below. Notice that
the rotating disk causes the pendulum to deviate from the sinusoidal-like behavior of the
undriven simple plane pendulum. We encourage you to try Problem 26 and explore how
θ(t) changes as you change the parameters in this example.
i i
i i
i i
Of course, this isn’t the end of the story when it comes to ODE solvers. Modern ODE
solvers often include an adaptive step-size control, where the value of dt changes. The value
of dt may increase when the function isn’t changing rapidly, and decreases when the function
does change rapidly. Equipped with adaptive step size, an ODE solver like RK4 can become
significantly more computationally efficient. Adaptive step-size controllers are beyond the
scope of this book, but the interested reader is encouraged to consult [Press et al. (2007)]
for more information.
It is often the case that you will use a numerical ODE solver that comes packaged
with your software of choice, or one that is included in a particular library. While you may
never need to write an ODE solver yourself, it is important to understand how they work,
so that you know how to interpret the results that they give. Errors from either misuse,
instability, or a variety of other mathematical reasons do occur. You need to be able to
identify when they happen and not assume the “computer is always right.” Mathematica’s
NDSolve is a powerful numerical ODE solver, which has built-in adaptive step size. The
NDSolve command has many features, including the ability to choose which algorithm
is being used. The interested reader should consult Mathematica’s online documentation
[MMA(2018)] for more information. Likewise, Python has the command odeint in the SciPy
library, which is based on the FORTRAN library ODEPACK. We comment on odeint in
more detail in Chapter 13. We have already seen both NDSolve and odeint used in this
book.
In the next section, we will return to the Lagrange formalism and study more closely
systems with constraints.
x2 + y 2 = `2 (8.6.1)
Notice that (8.6.1) takes the form,
i i
i i
i i
import numpy a s np
#d e f i n e c o n s t a n t s
a = 0.2 #r a d i u s o f d i s k
omega = 3/2∗ np . p i #a n g u l a r v e l o c i t y o f d i s k
g = 9.8 #a c c e l e r a t i o n due t o g r a v i t y
e l l = 1.0 #l e n g t h o f pendulum
d e f h ( time , a n g l e ) : #e q u a t i o n f o r ”v dot ”
r e t u r n −(a ∗omega ∗∗2/ e l l ) ∗np . c o s ( a n g l e −omega∗ time ) −(g / e l l ) ∗
,→ np . s i n ( a n g l e )
#d e i n e i n i t i a l c o n d i t i o n s
theta0 = 1.0 #i n i t i a l c o n d i t i o n f o r t h e t a
v0 = 0 #i n i t a l c o n d i t i o n f o r v
tmax = 2 0 . 0 #maximum time
dt = 0 . 0 1 #s t e p s i z e
n = i n t ( tmax/ dt ) #number o f s t e p s
t = np . l i n s p a c e ( 0 , tmax , n ) #a r r a y o f t i m e s
v = np . z e r o s ( n ) #a r r a y t h a t s t o r e s t h e s o l u t i o n
,→ f o r v
t h e t a = np . z e r o s ( n ) #a r r a y t h a t s t o r e s t h e s o l u t i o n
,→ f o r t h e t a
theta [ 0 ] = theta0
v [ 0 ] = v0
f o r i in range (1 , n) : #l o o p f o r f i n d i n g t h e s o l u t i o n
k1 = dt ∗ f ( v [ i −1])
j 1 = dt ∗h ( t [ i −1] , t h e t a [ i −1])
k2 = dt ∗ f ( v [ i −1]+ j 1 / 2 . 0 )
j 2 = dt ∗h ( t [ i −1]+dt / 2 , t h e t a [ i −1]+k1 / 2 . 0 )
k3 = dt ∗ f ( v [ i −1]+ j 2 / 2 . 0 )
j 3 = dt ∗h ( t [ i −1]+dt / 2 , t h e t a [ i −1]+k2 / 2 . 0 )
k4 = dt ∗ f ( v [ i −1]+ j 3 )
j 4 = dt ∗h ( t [ i −1]+dt , t h e t a [ i −1]+k3 )
t h e t a [ i ] = t h e t a [ i −1] + k1 / 6 . 0 + k2 / 3 . 0 + k3 / 3 . 0 + k4 / 6 . 0
v [ i ] = v [ i −1] + j 1 / 6 . 0 + j 2 / 3 . 0 + j 3 / 3 . 0 + j 4 / 6 . 0
i i
i i
i i
f (xi ; t) = 0 (8.6.3)
where the abbreviation xi says that f is a function of several coordinates (x, y, and z in
Cartesian coordinates) and the semicolon shows f ’s indirect dependence on time (through
the coordinates xi ). For example, we can rewrite (8.6.1) in the form of (8.6.2) by writing
f (x, y) = x2 + y 2 − `2 = 0. Constraints that must be written as relationships involving veloci-
ties are called nonholonomic constraints, unless the constraint equation can be integrated to
produce a relationship between coordinates. Although we won’t discuss nonholonomic con-
straints further, the interested reader can consult the classic text by Thornton and Marion
[Thornton and Marion(2004)].
Now we will derive Lagrange’s equations with undetermined multipliers. We will derive
the equations for a system with two degrees of freedom and then discuss a generalization for
systems with higher degrees of freedom. The system will have one holonomic constraint of
the form, f (x, y) = constant. We know that finding the equations of motion involves finding
functions x(t) and y(t) that make the action integral,
Z t2
S= (8.6.4)
L X, Ẋ, Y, Ẏ dt
t1
i i
i i
i i
Recall that ∂X/∂a = ηx and ∂Y /∂a = ηy . Solving (8.6.7) for ηy /ηx and inserting the result
into (8.6.6), we obtain:
Z t2
∂S ∂L d ∂L ∂L d ∂L ∂f /∂Y
= − − − ηx dt (8.6.8)
∂a t1 ∂X dt ∂ Ẋ ∂Y dt ∂ Ẏ ∂f /∂X
Setting ∂S/∂a = 0 at a = 0 gives (recall that X(a = 0, t) = x(t), and similar for Y (0, t)),
−1 −1
∂L d ∂L ∂f ∂L d ∂L ∂f
− = − (8.6.9)
∂x dt ∂ ẋ ∂x ∂y dt ∂ ẏ ∂y
The left-hand side of (8.6.9) depends on derivatives with respect to x and ẋ, while the
right-hand side depends on derivatives with respect to y and ẏ. Note that x and y are both
functions of t. In order for the two sides of (8.6.9) to be equal, they must both be equal
to a function of only t. We will define that function to be −λ(t), and we call λ Lagrange’s
undetermined multiplier. We do not know λ, hence the name “undetermined multiplier.”
Setting each side of (8.6.9) equal to −λ we obtain,
∂L d ∂L ∂f
− = − λ(t)
∂x dt ∂ ẋ ∂x
(8.6.10)
∂L d ∂L ∂f
− = − λ(t)
∂y dt ∂ ẏ ∂x
Equations (8.6.10) along with f (x, y) = constant will produce the equations of motion for
the system.
Of course, not all systems are limited to only two degrees of freedom and one constraint
equation. In general, systems will have n degrees of freedom and m equations of constraint.
We have already seen how to generalize the Lagrange formulation for s degrees of freedom,
in which case we have s Euler-Lagrange equations to solve. However, multiple constraint
equations means that there will be multiple derivatives dfj /da (j = 1, . . . , m, with m equal
to the number of constraints) that we will have to include in the derivation above. The
result is:
Euler-Lagrange Equation with Undetermined Multipliers
m
∂L d ∂L X ∂fj
− =− λj (t)
∂qi dt ∂ q˙i
j=1
∂qi (8.6.11)
fj (qi ; t) =0
i i
i i
i i
The second point is more important theoretically. It turns out that the Lagrange mul-
tiplier can tell us the force causing the constraint in the system. To illustrate this, we will
consider a simple system with two degrees of freedom and one constraint. The Lagrangian
is,
1
L = m ẋ2 + ẏ 2 − V (x, y)
2
and inserting this Lagrangian into the Euler-Lagrange equation yields,
∂V ∂f
− − mẍ = − λ
∂x ∂x
(8.6.12)
∂V ∂f
− − mÿ = − λ
∂y ∂y
We will focus on the equation involving x, which can be written in the form:
∂V ∂f
mẍ = − +λ (8.6.13)
∂x ∂x
The left-hand side mẍ is the x-component of the net force acting on the system. In
this case there are two forces acting on the system: an external force, associated with the
potential energy V (x, y), and the constraint force. The first term in the right-hand side of
(8.6.13) is the x-component of the external force acting on the system, gravity for example.
The remaining term λ∂f /∂x must therefore be the x-component of the constraint force.
We can repeat a similar argument for the y-equation (8.6.12) to see that λ∂f /∂y is the
y-component of the constraint force. In general we have,
where Fj,q
c
i
is the qi -component of the constraint force associated with the constraint equa-
tion fj (qi ; t) = constant. The following example illustrates the use of the Euler-Lagrange
equation to derive a constraint force.
θ R
i i
i i
i i
To begin the problem, we start by selecting polar coordinates to solve the problem.
Because the particle is moving on the surface of the sphere, our constraint equation is,
f (r; t) = r − R = 0 (8.6.15)
The force of constraint is the force that keeps the particle on the sphere and ṙ = 0.
When the particle leaves the sphere’s surface, ṙ 6= 0; therefore, there will be an acceleration
in the radial direction. We can find the relevant constraint force by finding:
∂f
Frc = λ
.
∂r
We know ∂f /∂r, so we need to find λ. We will find λ by setting up the Euler-Lagrange
equations for the system.
Using polar coordinates, we can write down the kinetic and potential energies for the
system:
1
T = m ṙ2 + r2 θ̇2
2
V =mgr cos θ
1
L = m ṙ2 + r2 θ̇2 − mgr cos θ
2
where we used the horizontal line in the above figure as y = 0 and therefore the potential
energy can be written as mgy with y = r cos θ. The Euler-Lagrange equations are:
∂L d ∂L ∂f
− =−λ (8.6.16)
∂r dt ∂ ṙ ∂r
∂L d ∂L ∂f
− =−λ . (8.6.17)
∂θ dt ∂ θ̇ ∂θ
mRθ̇2 − mg cos θ = − λ
(8.6.20)
−mgR sin θ − mR2 θ̈ =0
In the first equation above, we have found that λ = λ(θ, θ̇). The second equation will
give us a way of finding θ(t),
g
θ̈ = sin θ (8.6.21)
R
While we could attempt to integrate (8.6.21) and solve for θ, we actually don’t need
to do that. What we really need is θ̇, which can be found from (8.6.21) using,
i i
i i
i i
λ = mg (3 cos θ − 2)
2
θ0 = cos−1
3
Notice that this angle does not depend on the mass of the particle or the radius of the
sphere!
i i
i i
i i
s s
X ∂L X ∂L
dL = dqi + dq̇i = 0 (8.7.2)
∂qi ∂ q̇i
i=1 i=1
where s is equal to the particle’s number of degrees of freedom. Because dr is not a function
of time,
dqi d
dq̇i = d = (dqi ) = 0
dt dt
Therefore,
s
X ∂L
dL = dqi = 0 (8.7.3)
∂qi
i=1
i i
i i
i i
∂L
pi = = constant (8.7.6)
∂ q̇i
We have already seen that pi is the ith component of the particle’s generalized momen-
tum. The quantity pi is sometimes also referred to as the canonical momentum conjugate to
qi . From (8.7.4), we see that when L does not explicitly depend on the coordinate qi , then
the canonical momentum conjugate to qi is conserved. In other words, if qi doesn’t appear
in L, then pi is constant. Visual inspection of the Lagrangian can tell you what momenta are
conserved in the system! If qi doesn’t appear in L, then changing qi will not change L and
L is said to be “invariant” under changes of qi . Note that we can extend the work above to
N -particle systems by including an additional sum over α = 1, . . . , N . The procedure would
then follow as done above. Finally, when the Lagrangian is independent of the coordinate
qi , then qi is called a cyclic coordinate. Next, we will see that time-translation symmetry
in the Lagrangian will lead to conservation of energy.
For now, we will not restrict ourselves to Lagrangians that are time translationally invariant.
Next, we will insert the Euler-Lagrange equation
∂L d ∂L
= (8.7.8)
∂qi dt ∂ q̇i
into (8.7.7) to obtain,
s s
dL X d ∂L X ∂L ∂L
= q̇i + q¨i + (8.7.9)
dt dt ∂ q̇i ∂ q̇i ∂t
i=1 i=1
s
X d ∂L ∂L
= q̇i + (8.7.10)
dt ∂ q̇i ∂t
i=1
If the Lagrangian is invariant to time translations, then ∂L/∂t = 0 and therefore (8.7.10)
becomes,
s
" #
d X
L− pi q̇i = 0 (8.7.11)
dt
i=1
i i
i i
i i
where we used pi = ∂L/∂ q̇i . We see that we have a new conserved quantity in the square
brackets of (8.7.11). The conserved quantity is called the Hamiltonian function, or simply
the Hamiltonian. The Hamiltonian is defined as:
The Hamiltonian
s
X
H= pi q̇i − L (8.7.12)
i=1
Therefore !
X ∂rα X ∂rα
ṙ2α = ṙα · ṙα = q̇j · q̇k (8.7.15)
∂qj ∂qk
j k
i i
i i
i i
where
X ∂rα ∂rα
ajk = mα · (8.7.17)
α
∂qj ∂qk
and ajk = ajk (q1 , . . . , qn ). Next, we will compute the canonical momentum conjugate to qi ,
∂L ∂T X
pi = = = aij q̇j (8.7.18)
∂ q̇i ∂ q̇i
j
Proof of (8.7.18) is left as Problem 31. Finally, we insert (8.7.18) into (8.7.12):
X X
H=
aij q̇j q̇i − L
i j (8.7.19)
=2T − (T − V )
=T + V
i i
i i
i i
X
H= pi q̇i − L (qk , q̇k ; t) (8.8.3)
i
we find that the Hamiltonian is a function of generalized coordinates and generalized
momenta, H = H (qk , pk ; t).
Now we are ready to derive Hamilton’s equations of motion. We will derive Hamilton’s
equations for a system with one degree of freedom and leave the derivation for the s-degree of
freedom case as Problem 34. We begin by calculating the total differential of H = H(q, p; t),
∂H ∂H ∂H
dH = dq + dp + dt (8.8.4)
∂q ∂p ∂t
Next, we calculate dH using (8.8.3),
Hamilton’s equations of motion are also sometimes referred to as the canonical equations
of motion, and the description of a particle’s motion by these equations is referred to as
Hamiltonian dynamics. Regardless of the number of degrees of freedom, we also have the
relationship:
∂H ∂L
=− (8.8.11)
∂t ∂t
found by equating the coefficients of dt in (8.8.4) and (8.8.8).
Finally, if we divide (8.8.4) by dt, we find that,
dH ∂H
= (8.8.12)
dt ∂t
i i
i i
i i
Table 8.1: A summary of the Hamilton and Lagrange formulations of a system’s equations
of motion for a system with s degrees of freedom and no constraints.
While both the Lagrangian and the Hamiltonian can be used to find a system’s equa-
tions of motion, there are some differences between the two formulations. First, notice that
for a system of s degrees of freedom, Hamilton’s equations result in 2s first-order ODEs.
The Lagrange formulation for the same system results in s second-order ODEs. Hamilton’s
formulation used generalized coordinates and generalized momenta, while the Lagrange for-
mulation used generalized coordinates and generalized velocities. These different properties
are summarized in Table 8.1.
You now have three different methods of formulating a system’s equations of motion:
Newton’s second law, the Lagrangian, and the Hamiltonian. In addition, you can use con-
servation laws to set up and solve problems. The method you choose will depend on what
information you are given and what information you want to learn about the problem at
hand. In fact, you may learn different things by solving the same problem using multiple
techniques. Once students learn about Lagrangians, they sometimes ignore Newton’s sec-
ond law, but Newton’s second law can sometimes be the easiest one to use, especially when
dealing with drag forces and (as we will see later) noninertial reference frames. Use the right
tool for the job!
To end this section, we present a few examples of using the Hamilton formulation in
order to find equations of motion. Notice that all of these examples will have a common
theme: in order to find H, one needs to first find L and compute the generalized momenta.
y
θ ℓ
x
m
Solution:
From our work done earlier in this chapter, we know that the Lagrangian for this
system is (8.2.5),
1
L = m`2 θ̇2 + mg` cos θ
2
i i
i i
i i
where we used the coordinate θ instead of q. To find the Hamiltonian, we compute the
canonical momentum conjugate to θ,
∂L
pθ = = m`2 θ̇
∂ θ̇
Next, we use (8.7.12) to find the Hamiltonian,
H =pθ θ̇ − L
1 2 p θ 2
pθ
=pθ − m` + mg` cos θ (8.8.13)
m`2 2 m`2
p2θ
= − mg` cos θ
2m`2
Notice that we needed to substitute for θ̇ in order to get H in terms of θ and pθ .
Further, notice that the Hamiltonian is equal to the total mechanical energy of the system
because the transformation between Cartesian coordinates and polar coordinates is not
time-dependent.
Finally, we compute Hamilton’s equations,
∂H pθ
θ̇ = → θ̇ =
∂pθ m`2
(8.8.14)
∂H
ṗθ = − → ṗθ = −mg` sin θ
∂θ
There are a few things to notice with our solution. First, Hamilton’s equation for θ̇
simply reproduces what was obtained by computing the canonical momentum. Second,
the canonical momentum in this case is actually the angular momentum, not the linear
momentum. Finally, notice that ṗθ is the torque acting on the system.
We plot a numerical solution to (8.8.14) using the following Mathematica code. Here
we used m = 1.0 kg, ` = 1.0 m, and g = 9.8 m/s2 with initial conditions θ(0) = 3.0 rad and
pθ (0) = 0.
l = 1.0;
m = 1.0;
g = 9.8;
soln = NDSolve[{θ0 [t] == pθ[t]/(m ∗ l∧ 2), pθ0 [t] == −m ∗ g ∗ l ∗ Sin[θ[t]], θ[0] == 3.0,
pθ[0] == 0}, {θ, pθ}, {t, 0, 20}];
Plot[θ[t]/.soln, {t, 0, 20}, BaseStyle → {FontSize → 18}, Frame → True, Axes → False,
FrameLabel → {“time”, “θ”}, ImageSize → Large]
The resulting plot is shown below.
i i
i i
i i
θ
Out[ ]=
-1
-2
-3
0 5 10 15 20
time
ρ=R
m
r
z
x2 + y 2 = R 2 (8.8.15)
which can also be written as ρ2 = R2 . Therefore, there are two degrees of freedom in this
problem.
Next, we compute the kinetic energy. In cylindrical coordinates, v 2 = ρ̇2 + ρ2 θ̇2 + ż 2 .
However, using our constraint ρ̇ = 0, the kinetic energy of the particle becomes,
1
T = m R2 θ̇2 + ż 2 (8.8.16)
2
The potential energy is best found in Cartesian coordinates first. We are given that
F = −kr = −krr̂ = −k xî + y ĵ + z k̂ (8.8.17)
i i
i i
i i
1
V = kr2
2 (8.8.18)
1
= k x2 + y 2 + z 2
2
If we apply the constraint (8.8.15), then the potential energy becomes
1
V = k R2 + z 2
2
We next compute the Lagrangian L = T − V ,
1 1
L = m R2 θ̇2 + ż 2 − k R2 + z 2 (8.8.19)
2 2
From the Lagrangian we see that the canonical momentum conjugate to θ is conserved
because θ is a cyclic coordinate. In order to find the Hamiltonian, we need to compute pθ
and pz ,
∂L ∂L
pθ = = mR2 θ̇ and pz = = mż (8.8.20)
∂ θ̇ ∂ ż
Note that pθ again is an angular momentum, whereas pz is the translational momen-
tum. Next, we use (8.8.3) to find the Hamiltonian,
H =pθ θ̇ + pz ż − L (8.8.21)
1 p 2 p 2 1
pθ pz θ z
=pθ + 2
+ 2
+ 2
(8.8.22)
pz − m R − k R z
mR2 m 2 mR2 m 2
p2
p 2 1
= θ 2 + z + k R2 + z 2 (8.8.23)
2mR 2m 2
Finally, we use the Hamiltonian to compute Hamilton’s equations,
∂H pz
ż = → ż =
∂pz m
∂H pθ
θ̇ = → θ̇ = 2
∂pθ mR (8.8.24)
∂H
ṗz = − → ṗz = − kz
∂z
∂H
ṗθ = − → ṗθ =0
∂θ
Notice that ṗz is equal to the z-component of the external force acting on the system.
i i
i i
i i
constrained to move along the surface of a cylinder. Recall that θ was a cyclic coordinate
in Example 8.9. The equations of motion relevant to θ were,
pθ
θ̇ = and ṗθ = 0.
mR2
Because pθ is constant, it is easy to find θ(t),
Z t
pθ pθ
θ(t) = 2
dt0 = 2
(t − t0 ) .
t0 mR mR
When the canonical momentum is constant, the equation of motion can be reduced to
quadrature (as above). If a coordinate system can be found such that all coordinates are
cyclic, then the equations of motion are trivial. In fact, it is possible to find such coordinate
systems. Those transformations were developed by Carl Gustov Jacob Jacobi (1804–1851)
and the work is referred to as the Hamilton-Jacobi theory. Hamilton-Jacobi theory is beyond
the scope of this book, but the motivated reader should consult [Goldstein et al. (2001)].
Furthermore, simply knowing the Hamiltonian alone can give us insight into all possible
behaviors of the system. This can be done by using the Hamiltonian to plot the system’s
phase space. Unlike in Chapter 6 where we plotted x versus v, we will plot p versus q
to create the phase space. In particular, we will plot the canonical momentum versus its
conjugate generalized coordinate. In general, the dimension of the phase space is 2s, where
s is the system’s number of degrees of freedom. Notice that H = H (qk , pk ) and will define
surfaces in the phase space. To understand the value of such plots, we will revisit the simple
pendulum, yet again.
p2
H(q, p) = − mg` cos q (8.9.1)
2m`2
where we have defined p = pθ and q = θ for simplicity of notation. The Mathematica code
to create a contour plot with m = 1.0 kg, ` = 1.0 m, and g = 9.8 m/s2 is shown below.
l = 1.0;
m = 1.0;
g = 9.8;
2
p
H[q , p ]:= 2∗m∗l 2 − m ∗ g ∗ l ∗ Cos[q];
ContourPlot[H[q, p], {q, −π, π}, {p, −10, 10}, ContourShading → None,
Contours → {−m ∗ g ∗ l ∗ Cos[0.5], −m ∗ g ∗ l ∗ Cos[2], m ∗ g ∗ l, 2 ∗ m ∗ g ∗ l, 3 ∗ m ∗ g ∗ l},
ContourStyle → {Black, Black, {Red, Dashed}, Black, Black}]
The resulting plot is shown as follows.
i i
i i
i i
e
d
c
b a
d
e
The Hamilton formulation is a powerful one, especially for theoretical mechanics. Exten-
sions of the Hamilton formulation like the Hamilton-Jacobi theory laid the groundwork for
theoretical structure of quantum mechanics. Beyond [Goldstein et al. (2001)], the very moti-
vated student will find [Arnold(1997)] to be both a highly challenging and highly rewarding
presentation of theoretical mechanics that goes well beyond the undergraduate-level, tying
classical mechanics to advanced mathematics such as Lie algebra and symplectic geometry.
i i
i i
i i
stationary. The functions qj (t) that make the action integral stationary can be found using
the Euler-Lagrange equation,
d ∂L ∂L
− =0
dt ∂ q̇i ∂qi
The Euler-Lagrange equation produces s second-order ODEs, where s is the number
of degrees of freedom in the system, which can be solved for qi (t).
In the case, where there are m holonomic constraints on the system, we can express
the constraints as m constraint equations fj and the Euler-Lagrange equation becomes
m
∂L d ∂L X ∂fj
− =− λj (t)
∂qi dt ∂ q˙i ∂qi
j=1
In addition to using the Lagrangian, the Hamiltonian can also be used to derive a
system’s equations of motion. The Hamiltonian is defined as,
s
X
H= pi q̇i − L.
i=1
i i
i i
i i
projectile travels a short enough distance that the Coriolis force (see Chapter 10) is
negligible.
5. Consider a disk of radius R rolling down an inclined plane. What quantities are needed
to describe the state of the system? Are there constraints among those quantities?
Hint: The disk is rolling, without sliding.
6. A very small ball of mass m and radius a is constrained to roll along the surface of a
sphere of R a. The ball is constrained to move along geodesics on the sphere. How
many degrees of freedom does the ball have?
Section 8.3: Hamilton’s Principle
7. Consider a single-particle system of s degrees of freedom. The particle experiences a
force which has an associated potential energy V (qi ). Show that the functions qi (t)
Rt
for i = 1, . . . , s that make the action integral S = t12 L (qi , q̇i ) dt stationary also obey
Newton’s laws. In other words, extend the work done in Section 8.3 to a single-particle
system with s-degrees of freedom.
8. Repeat Problem 7 for an N -particle system.
Section 8.4 and 8.6: Examples of Lagrangian Dynamics
Note that the problems in this section combine problems for sections 8.4 and 8.6.
9. Fill in the necessary steps to calculate (8.4.19) and (8.4.20) in Example 8.3.
10. In Example 8.5, find the acceleration of the inclined plane and the particle in the limit
of M → 0. Discuss your results.
11. A simple plane pendulum made of a mass m and massless rod of length ` is attached
to a support that accelerates upward with an acceleration a. Find the equations of
motion and the period of small oscillations.
12. Find the equation of motion for a simple plane pendulum made of a mass m and a
massless rod of length ` which is attached to a support that oscillates horizontally,
such that xs = a cos (ωt), where xs is the location of the support on the horizontal
axis. The support has no vertical motion.
13. A simple plane pendulum made of a mass m and massless rod of length ` is attached to
a support that oscillates vertically, such that ys = a cos (ωt) , where ys is the location
of the support on the vertical axis. The support has no horizontal motion. Find the
equation of motion and period of small oscillations.
14. A double pendulum consists of two pendulums attached in series as shown in Figure
8.5 that follows. Find the equation of motion for the double pendulum.
15. Consider a simple pendulum of mass m and initial length `. After the pendulum
begins to swing, its length is shortened at a rate `˙ = −α, where α is a constant.
Use the Lagrangian to compute the equations of motion of the system. Is the energy
conserved in this system?
16. A small particle with mass m slides down a smooth circular wedge with mass M and
radius of curvature R as shown in Figure 8.6. The wedge is free to move horizontally
along a frictionless surface. Find the equation of motion for each object.
i i
i i
i i
ℓ1
θ1
m1
ℓ2
θ
m2
17. Find the equations of motion for a single particle in three dimensions experiencing a
central force field. Use spherical coordinates.
18. A particle of mass m is constrained to move along a frictionless massless circular hoop
of radius a. The hoop rotates about its vertical diameter with a constant angular
speed ω. Find the equilibrium positions of the particle.
19. A massless spring of equilibrium length ` connects two masses, m1 and m2 . The masses
are free to rotate and oscillate on a horizontal frictionless surface. Find the equations
of motion for the particles.
20. Consider a mass attached to a massless spring with a spring constant k and equilib-
rium length ` that is allowed to swing on a frictionless plane. This is the so-called
elastic pendulum. Find the equation of motion for an elastic pendulum of mass m and
equilibrium length `.
21. A smooth horizontal circular wire of radius R rotates about a point on its perimeter. A
particle of mass m is constrained to move along the wire as shown in Figure 8.7. Find
the equation of motion for the particle. About what point does the particle oscillate?
i i
i i
i i
ω
x
22. A rope of mass M and length L is stretched out along a frictionless table such that
a length z0 is hanging over the edge of the table. Find the equation of motion of the
rope if it is released from rest. Show that the rope requires a time
s
L
t= cosh−1 (L/z0 )
g
to completely slide off the table. You may remember this problem from Chapter 2,
where we did it using Newton’s laws.
23. A particle of mass m moves in a plane under the force F = −αr2 where F points
towards the origin. Find the equations of motion for the particle. Are there any con-
served quantities? Is so, what are they?
24. A sphere of radius a is constrained to roll without slipping on the inside of a cylinder
with a radius R. Find the equations of motion of the sphere and calculate the period
of small oscillations.
25. A particle of mass m is constrained to move on a frictionless wire is bent into the
shape of a helix with r = R and z = αθ, where R and α are positive constants. The
particle is under the influence of gravity, which points in the −z direction. Find the
equation of motion for the particle. What is the particle’s acceleration when R → 0?
Section 8.5: Numerical Solutions to ODE’s Using the Fourth-Order Runge-
Kutta Method
26. Repeat Example 8.6 for different values of the parameters R, ω, and `. Discuss the
different types of motions that result from changing the parameters.
27. Write a program that uses RK4 to numerically solve the equation of motion for the
simple harmonic oscillator, ẍ + ω02 x = 0 where ω0 = 1.0 rad/s with initial conditions
x(0) = 1.0 m and ẋ(0) = 0.5 m/s.
28. Using the same equation and values for the SHO in Problem 27, calcu-
late the error in x(t) at t = 3.0 s for the RK4 and Euler method by using
dt = 0.001, 0.01, 0.1, 0.2, and 1.0. We will define the error as,
exact value − result
exact value
where “exact value” is calculated from the closed-form solution of the SHO and
“result” is the result obtained by either the RK4 or Euler method.
i i
i i
i i
29. Use an RK4 program to solve the damped driven harmonic oscillator,
for ω = 1.0 rad/s, β = 2.0 rad/s, ω0 = 0.5 rad/s, and A = 0.25 m/s2 with initial con-
ditions x(0) = 1.0 m and ẋ(0) = 0 m/s.
ẋ =σ(y − x)
ẏ =rx − y − xz
ż =xy − bz
are an important set of equations in the field of nonlinear dynamics. Using an RK4
program, solve for x(t), y(t), and z(t) for parameter values: σ = 10, r = 28, and b = 8/3
and initial conditions x(0) = 0.1, y(0) = z(0) = 0. Plot each solution individually and
then make a three-dimensional plot which contains x(t) vs. y(t) vs. z(t).
Section 8.7: Conservation Theorems and the Lagrangian
31. Prove (8.7.18):
∂L ∂T X
pi = = = aij q̇j
∂ q̇i ∂ q̇i
j
32. A particle of mass m moves within a force field that has an associated potential energy
V = −k/r where k is a positive constant. What are the conserved quantities for this
particle?
33. Consider a system with s degrees of freedom and Lagrangian L(qk , q̇k ), where k =
1, . . . , s. Now, let F = F(qk ) be any function of the generalized coordinates, qk . We
can construct an additional Lagrangian,
L0 = L + dF/dt
i i
i i
i i
39. Find Hamilton’s equations for a particle of mass m confined to the surface of a cone,
which in cylindrical coordinates is described by r = az, inside a uniform gravitational
field. Show that for a given energy, there are maximum and minimum values of z for
which the particle is confined.
40. Consider the force
k
F = − e−λt (8.11.1)
x
acting on a particle of mass m moving in one dimension. Assume that the constants
k and λ are bothR positive. We can create a “pseudo potential energy,” using the rela-
tionship, V = − F dx. Using the pseudo potential energy find Hamilton’s equations
of motion for the particle. Compute ∂H/∂t and ∂L/∂t and show that they satisfy
(8.8.11):
∂H ∂L
=−
∂t ∂t
41. Recall the cycloid from the brachistochrone problem,
y =a (1 − cos θ)
x =a (θ − sin θ)
i i
i i
i i
44. Both in Chapter 5 and in Problem 36 of this chapter, we studied the double-well
potential, V = − 12 kx2 + 14 x4 . If you haven’t already done so, find the Hamiltonian
for this potential. Using the Hamiltonian, plot the phase space for = 1 and various
values of k. How does the phase space change? For each value of k that you try,
describe the possible motions of the particle. Draw the phase space plots and identify
the types of motion associated with each region of the phase space.
i i
i i
i i
CHAPTER 9
In this chapter, we will first look at the general concept of central forces and their importance
in describing physical phenomena. We will begin with a general discussion of central forces.
Next we will examine the problem of two objects interacting via a central force, this is the
so called two-body problem. Our analysis of the two-body problem will lead to Kepler’s
laws of planetary motion. We will then examine the specific case of a planet orbiting a
star. Finally, we will study the three-body problem where three objects interact via central
forces. Solving for the motion of one of the objects in the three-body problem will require
us to use the computational skills developed so far in this book.
F = f (r)r̂ (9.1.1)
where f (r) is a scalar function of the magnitude r of the position vector r, and r̂ is a unit
vector along the position vector as shown in Figure 9.1.
273
i i
i i
i i
m1 r
F21
r
F12
m2
Figure 9.1: Two masses m1 and m2 interacting via a central force F.
The best-known examples of central forces in physics are Newton’s law of universal
gravitation FG (r) between two masses m1 and m2 , and the electrostatic Coulomb force
FC (r) between two electric changes q1 and q2 . The mathematical description of these two
forces is very similar:
Gm1 m2
FG (r) = − r̂ (9.1.2)
r2
1 q1 q2
FC (r) = r̂ (9.1.3)
4π0 r2
where G = 6.67×10−11 m3 kg−1 s−2 is the universal gravitational constant and 0 = 8.85 ×
10−12 m−3 kg−1 s4 A2 is the corresponding constant for electrical forces, called the permit-
tivity of free space.
1 1 1 ∂Fr
∂ ∂Fθ ∂
∇×F = Fφ sin θ − r̂ +
− rFφ θ̂
r sin θ ∂θ ∂φ r sin θ ∂φ ∂r
1 ∂ (rFθ ) ∂Fr
+ − φ̂ (9.1.4)
r ∂r ∂θ
∇×F = 0 (9.1.5)
and therefore, F is a conservative force.
As we saw in Chapter 5, conservative forces can be associated with potential energy
functions V (r) = − F · dr. Hence, central forces have a potential energy associated with
R
them. The corresponding scalar gravitational potential energy VG and electrostatic potential
i i
i i
i i
energy VC are:
Gm1 m2
VG (r) = − (9.1.6)
r
1 q1 q2
VC (r) = (9.1.7)
4π0 r
Another well-known example of a central force is the intermolecular long-range van der
Waals force, which is of importance in several branches of science. In one of its simplest
forms, the van der Waals force between a pair of neutral atoms or molecules can be approx-
imated by the Lennard-Jones potential (also termed the L-J potential or 6-12 potential).
The mathematical form of the 6-12 potential VLJ (r) is:
hc c2 i
1
VLJ (r) = 12 − 6 (9.1.8)
r r
where c1 and c2 are constants with the appropriate SI units.
In the field of particle and atomic physics, the Yukawa short-range nuclear force is of
great importance, and the corresponding Yukawa potential VY (r) is represented mathemat-
ically by:
e−c5 r
VY (r) = c4 (9.1.9)
r
where c4 , and c5 are constants with the appropriate SI units.
where Fi is the force acting on the particle with mass mi . In Chapter 4, we further simplified
this equation to show,
N
X N X
X
N= ri × Fext
i + (ri − rj ) × Fij (9.1.11)
i=1 i=1 j>i
where Fext is the sum of the external forces acting on mi and Fij is the force of interaction
between mi and mj . If we assume an isolated system then Fext = 0 and the first term
in (9.1.11) is equal to zero. Using a coordinate system similar to that in Figure 9.1, we
can write ri − rj = rij r̂ij , where r̂ij is a unit vector that lies along a line joining mi and
mj . Furthermore, if the force of interaction between the particles is a central force, then
Fij = Fij (r)r̂ij and the second term in (9.1.11) is: (rij r̂ij ) × (Fij r̂ij ) = 0. Therefore the net
torque acting on the system is zero, and angular momentum is conserved.
In the next section, we will examine the so-called two-body problem in detail. This
discussion will lead us to Kepler’s famous laws of planetary motion.
i i
i i
i i
r1 r=r1-r2
R
m2
r2
O
Figure 9.2: The position vectors r1 and r2 and relative distance r = r1 − r2 of the two masses
interacting with a central force F (r)r̂. The white dot at the location R shows the location
of the system’s center of mass.
If r1 and r2 are the position vectors of the two masses m1 and m2 , respectively, then the
vector distance between the two masses is
r = r1 − r2 (9.2.1)
1 1
r̈1 −r̈2 = − + F (r)r̂ (9.2.4)
m1 m2
m1 m2
(r̈1 −r̈2 ) = −F (r)r̂ (9.2.5)
m1 + m2
which can be rewritten as
Newton’s Second Law for the Reduced Mass µ
i i
i i
i i
This equation tells us that the two-body system behaves just like a single mass µ with a
position vector r, moving under the influence of the central force F = F (r)r̂. This single
mass µ is called the reduced mass of the two-body system and represents an effective inertial
mass appearing in the two-body problem. Using the reduced mass µ allows us to study the
two-body problem as if it were a one-body problem. Note that the minus sign appears in
(9.2.6) because we assumed an attractive force. A repulsive force would give similar results
but without the minus sign in (9.2.6).
What is the reduced mass? The definition in (9.2.7) can be difficult to interpret. Consider
the Earth-Sun system. In this case the mass of the Sun mS is much greater than the mass
of the Earth, mE . If we use m1 = mS and m2 = mE in (9.2.7), we find that µ ≈ mE . Hence,
when we study the motion of the Earth-Sun system, we are essentially studying the motion
of the Earth. As the masses become more comparable, we cannot associate the reduced
mass with either one of the masses, as we will demonstrate later in Example 9.1.
When working with a system of multiple particles, we know that the system’s center of
mass can be useful in describing the behavior of the system. As we discussed in Chapter 4,
the position of the center of mass of the two-body system is defined by:
m1 r1 + m2 r2
R= (9.2.8)
m1 + m2
Because the translational motion of the system as a whole is not of interest (we are
focused on the particles’ orbits with respect to each other), we can set R = 0, as shown
in Figure 9.3a. In other words, we can place the center of mass at the origin. This new
coordinate system is different from that shown in Figure 9.2. As we will see, this new
coordinate system with the center of mass at the origin, will be very useful in describing
the orbits.
Figure 9.3: (a) The binary system can be best described using a system of coordinates axes
with the origin located at the center of mass of the system. (b) The motion can also be
described by Newton’s law (9.2.6), as the motion of the reduced mass µ, with position vector
r = r1 − r2 , moving with the relative velocity of the two masses v = ṙ1 − ṙ2 .
Using R = 0, we have,
m1 r1 + m2 r2 = 0 (9.2.9)
By combining (9.2.9) with the definition r = r1 − r2 for the position vectors r1 and r2 , we
obtain:
m2
r1 = r (9.2.10)
m1 + m2
i i
i i
i i
m1
r2 = − r (9.2.11)
m1 + m2
The equations (9.2.10) and (9.2.11) allow us to relate the particles’ position to the
separation vector r. In addition, we can write the Lagrangian for the system as:
1 1
L = m1 ṙ21 + m2 ṙ22 − V (r) (9.2.12)
2 2
where V (r) = − F (r)dr is the potential energy associated with the central force F. Sub-
R
stituting the values of (9.2.10) and (9.2.11) in the Lagrangian (9.2.12), we find after some
simple algebra that the Lagrangian of the two-body system is:
Similar to Newton’s second law, the Lagrangian shows that we can reduce the two-body
problem to the motion of a single object with a mass equal to the reduced mass µ, moving
with the relative velocity of the two masses v = ṙ1 − ṙ2 .
Similarly, the total angular momentum of the two-body problem is found from:
dr1 dr2
L = m1 r1 × + m2 r2 × (9.2.14)
dt dt
Substituting (9.2.10) and (9.2.11) into (9.2.14) we can obtain (see Problem 10 at the end
of this chapter):
` = µ (r × ṙ) (9.2.15)
where we have replaced L (for two bodies) with ` (for one body). This equation tells us that
the total angular momentum of the two-body problem is equal to the angular momentum of
the reduced mass µ, located at the point r = r1 −r2 and moving with velocity v = ṙ = ṙ1 − ṙ2 .
Again, the result is that we can reduce the two-body problem to studying the motion of a
single particle of mass µ.
Furthermore, taking the time derivative of (9.2.15) results in `˙ = 0, hence just as before,
the angular momentum is conserved. One of the physical consequences of the conservation
of the angular momentum vector ` is that the motion of the mass µ has to remain on
the same plane (which also contains the force center), in order to keep the direction of `
constant at all times. This is shown for the special case of a circular orbit in Figure 9.4.
Because the motion of the particle is restricted to a plane, polar coordinates are a useful
coordinate system when dealing with central force problems. The constant nature of ` helps
explain why all of the orbits of the planets in our solar system are essentially co-planar and
continue to stay that way.
i i
i i
i i
=r×p
θ r
p
μ
Figure 9.4: Relationship of the angular momentum vector `, position vector r and linear
momentum vector p = µv.
Next, we will compute the magnitude of the angular momentum in polar coordinates.
Recalling that in polar coordinates, ṙ = ṙr̂ + rθ̇θ̂θ , (9.2.15) yields:
Magnitude of the Angular Momentum of the Two-Body Problem
` = µr2 θ̇ (9.2.16)
We can also compute the total energy of the two-body system using E = T + V ,
1
E = µ ṙ2 + r2 θ̇2 + V (r) (9.2.17)
2
2
1 2 1 2
`
= µṙ + µr + V (r) (9.2.18)
2 2 µr2
which can be simplified to
Total Energy of the Two-Body System
2
`2 1 dr
E= + µ + V (r) (9.2.19)
2µr 2 2 dt
Again, the physical interpretation of this equation, and the others before it, is that the two-
body system can be replaced by a single reduced mass µ at a distance r =| r1 − r2 |from the
origin (i.e., the force center and the system’s center of mass), and moving with the velocity
v = ṙ = ṙ1 − ṙ2 . This is shown schematically in Figure 9.3(b).
While the reduced mass is a useful tool for performing the calculations above, what we
really want is to be able to calculate the motions of m1 and m2 . In the next section, we will
derive the equations of motion for the two-body problem.
i i
i i
i i
which, in principle, can be integrated to get t(r) and then inverted to get r(t), once we have
specified V (r). However, r(θ) would be a more useful equation of motion because it would
give us a mathematical formula for the observed motion. We can substitute,
dr dr dθ dr `2
= = (9.3.2)
dt dθ dt dθ µr2
into (9.3.1) to obtain
Z
`dr
θ(r) = ± r (9.3.3)
`2
r2 2µ E − V (r) − 2µr 2
In principle, after specifying V (r), we could perform the integral to get a mathematical
form for θ(r), and then invert it to get r(θ). However, in general, the integral in (9.3.3) is
difficult to do, and the inversion is also not easy.
Computers can provide a means of obtaining numerical solutions, in which case it is often
easier working directly with Newton’s second law. In the next two examples, we show how
numerical solutions for an object’s orbit can be obtained by numerically solving Newton’s
second law.
Example 9.1: The orbits in a two-body gravitational system
Integrate Newton’s second law using Newton’s law of universal gravitation to represent
for the interaction force between two masses m1 and m2 , in order to obtain the orbits of
a binary star system with masses m1 = 1, m2 = 2 (in arbitrary units). In order to simplify
the presentation of the numerical results, use arbitrary units with G = 1. At time t = 0 the
reduced mass is located at x = 1, y = 0 and has a velocity with an x-component of vx = 0
and a y-component of vy = 1.
Solution:
As discussed in this section, Newton’s second law for the two masses becomes a single
equation for the reduced mass:
Gm1 m2
µr̈ = − r̂
r2
where we have used Newton’s law of universal gravitation as the central force. Now we
need to rewrite the above formula in terms of Cartesian coordinates. The unit vector is
given in Cartesian coordinates (x, y) by:
r xî + y ĵ
r̂ = =p
r x2 + y 2
d2 x Gm1 m2 d2 y Gm1 m2
µ =− x µ =− y (9.3.4)
dt2 (x + y 2 )
2 3 /2 dt2 (x + y 2 )
2 3 /2
i i
i i
i i
In the Mathematica code below, the command NDSolve is used to solve the system of
differential equations (9.3.4), and the Table command is used to construct a list of the
pairs of position coordinates r = (x, y) from t = 0 to t = 10, in time steps of dt = 0.01. The
coordinate pairs are stored in the variable xyList. We see that the reduced mass µ follows
an elliptical orbit. As we will see later, this is an example of Kepler’s First Law.
We also plotted the positions of the two masses. The positions can be calculated from
r by using (9.2.10) and (9.2.11); this is done easily by multiplying the list xyList with the
factors m1 /(m1 + m2 ) and −m2 /(m1 + m2 ), to get r1 and r2 , respectively. Finally a graph
of the orbits of the two masses is produced by using the ListPlot command. Notice that
each mass follows an elliptical path and the center of mass (at the origin) is one of the
foci of each ellipse.
m1 = 1; m2 = 2; x0 = 1; y0 = 0; vx0 = 0; vy0 = 1;
µ = m1 ∗ m2/(m1 + m2);
x1 = x/.sol[[1]];
y1 = y/.sol[[1]];
xyList = Table[{x1[t], y1[t]}, {t, 0, 10, .01}];
OUTPUT:
location of center of ellipse=0.400014
location of CM=0.0000323403
GraphicsGrid[
{{Show[ListPlot[xyList], Graphics[Disk[{centerofMass, 0}, .02]]],
ListPlot[{−m2/(m1 + m2) ∗ xyList, m1/(m1 + m2) ∗ xyList}]}}, ImageSize->Large]
i i
i i
i i
0.3
0.4
0.2
0.2 0.1
-0.2 0.2 0.4 0.6 0.8 1.0 -0.6 -0.4 -0.2 0.2
-0.2 -0.1
-0.2
-0.4
-0.3
Although numerical solutions are very useful, it would be helpful to have an analytical
solution. Next, we derive a general equation, which provides another means of computing the
reduced mass’ orbit for the two-body problem. It can also be used to find the mathematical
form of the central force acting on the reduced mass, if the mathematical form of the orbit
is known to us. In order to obtain the central force, we will need to use Newton’s second
law in polar coordinates. Here r and θ are the usual position and angle polar coordinates.
The acceleration in polar coordinates is given by:
dv
= r̈ − rθ̇2 r̂ + rθ̈ + 2ṙθ̇ θ̂ (9.3.5)
dt
Newton’s second law becomes:
i i
i i
i i
or
d2 u µ
+ u = − 2 2 f u−1 (9.3.14)
dθ 2 ` u
Finally this equation can be written in terms of r by substituting u = 1/r:
d2 1 1 µr2
2
+ = − 2 f (r) (9.3.15)
dθ r r `
This equation is useful when we know the orbit in the form r(θ), and we want to evaluate
the central force f (r) that creates this orbit.
Example 9.2: Finding the central force when we know the orbit
The orbit r(θ) of a mass m moving inside a central force field, is given by the expression:
r(θ) = kθ
`2 1 1
2
d
f (r) = − 2 +
µr dθ2 r r
We now have
d2 1 d2 1 2
= 2 = 3
dθ2 r dθ kθ kθ
and the f (r) equation above gives:
`2 2k 2 1 `2 2k 2 1
f (r) = − 2 + =− + 3
µr r3 r µ r5 r
So the force in this system is a linear combination of two terms 1/r5 and 1/r3. .
(b) We can find the orbit θ(t) by using the conservation of angular momentum:
dθ
µr2 θ̇ = µr2 =`
dt
and substituting r = kθ:
dθ µk 2 d (θ)3
µk 2 θ2 = =`
dt 3 dt
This can be integrated to yield:
3t`
θ3 (t) = +C
µk 2
i i
i i
i i
and using the initial condition θ(0) = 0 yields the orbit θ(t), or equivalently the orbit
r(t) = kθ(t): s s !
3t` 3 3t`
θ(t) = 3
r(t) = k
µk 2 µk 2
Example 9.3: Finding the total energy when we know the orbit
Find the total energy E for the orbit r(θ) of the mass µ in Example 9.2.
Solution:
We find the potential energy V (r) = − f (r) dr by using the expression for the force
R
` 2k 1 1
Z 2 2
`2 k 2
Z
V (r) = − f (r) dr = + dr = − +
µ r5 r3 2µ r4 r2
We can evaluate dr/dt using the chain rule and conservation of angular momentum ` =
µr2 θ̇:
dr dr dθ d (kθ) dθ `
= = = k θ̇ = k 2
dt dθ dt dθ dt µr
Finally the total energy will be:
2 2
`2 1 `2 1 `2 k 2 1
dr `
E = T +V = + µ + V (r) = + µ k 2 − + =0
2µr2 2 dt 2µr2 2 µr 2µ r4 r2
w00 = −w (9.4.3)
Notice that the result is a simple harmonic oscillator equation for w with a natural frequency
ω0 = 1. The solution is w = A cos θ + B sin θ. With a careful choice of coordinates, we can set
i i
i i
i i
B = 0. Let’s choose θ = 0 as the angular position of µ when it makes its closest approach
rmin to the force center (located at the center of mass). At this point, ṙ = 0 (because this
is a turning point in the orbit) and due to (9.3.11), u0 = 0. Hence, we are setting an initial
condition for u to be u(θ = 0) = 0. Because w0 = u0 then w0 (θ = 0) = 0 and we can choose
B = 0. Hence we have,
µk
u = A cos θ + (9.4.4)
`2
If we set e = `2 A/µk and α = `2 /µk, we get
1 1
= (1 + e cos θ) (9.4.5)
r α
which is the equation of a conic section in polar coordinates with one focus at the origin,
e is called the eccentricity and α is called the semi-latus rectum. Later in this chapter, we
will discuss the equation for conic sections in polar coordinates in more detail.
When studying planetary motion, Johannes Kepler (1571–1630) developed three laws of
planetary motion. Kepler’s first law stated that planets orbit in an ellipse and that the Sun
is at one of the foci of the ellipse. An ellipse, of course, is one example of a conic section.
However, starting from Newton’s second law, we have found a much more general result.
We have found that the motion of the reduced mass is of the form of a conic section: circle,
ellipse, parabola, or hyperbola. What conditions dictate which conic section the reduced
mass follows? In the next section, we will begin to answer this question.
However, before we move on, there is one piece of (9.4.5) that is still unsatisfactory, the
eccentricity e is in terms of the constant of integration A. It would be better if we could
rewrite e in terms of physical parameters, similar to α. We begin by using our choice of
coordinates where ṙ = 0 at r = rmin and rmin occurs at θ = 0, hence (9.4.5) gives
α
rmin = (9.4.6)
1+e
and the total energy at rmin is
`2 k
E= − (9.4.7)
2µrmin
2 rmin
1 l2
E = T + V = µṙ2 + + V (r) (9.5.1)
2 2µr2
i i
i i
i i
where we have returned to the generic potential energy. Note that for an inverse square law,
V = −k/r. The first term in (9.5.1) can be interpreted as a kinetic energy term, because it
is of the form of one half times a mass and the square of the radial velocity. It is convenient
to think of the second two terms as potential energy terms.
We notice that the term
l2 1
VC ≡ = µr2 θ̇2 (9.5.2)
2µr2 2
can be interpreted as a form of energy called the centrifugal potential energy VC , and the
corresponding force will be the centrifugal force FC :
∂VC l2
FC = − = 3 (9.5.3)
∂r µr
Note that FC = FC r̂ points radially outward from the center of mass (the force center for
the two-body problem). We can then write the effective potential energy Veff (r) as the sum
of two terms, the centrifugal potential VC and the potential energy V (r) corresponding to
the central force.
The Effective Potential in a Central Force Field
l2
Veff (r) = + V (r) = VC + V (r) (9.5.4)
2µr2
Figure 9.5 shows a plot of the effective potential Veff (r) as a function of the distance r, and for
an attractive potential V = −k1 /r where k1 = 2.5 in SI-units. The dashed lines in this graph
indicates the repulsive positive centrifugal potential VC = 1/r2 , and the attractive potential
V = −k1 /r. The thick solid line in this figure shows the sum of these two potentials, which
represents the effective potential Veff (r). When the distance r → 0 the effective potential
goes to infinity, while as r → ∞ the effective potential goes to zero. The horizontal lines in
this figure indicate two possible energies of the reduced mass, E0 = −1.55 J and E1 = −1.3 J.
Note that the value of these energies are arbitrarily chosen for illustrative purposes below.
We can use the methods of Chapter 5 in order to qualitatively describe the motion of
the reduced mass µ as a function of its total energy E. Recall, that we can think of the
reduced mass as a particle that is rolling along the track made by the shape of the Veff
graph. When E = E0 , the particle is fixed at the minimum of Veff . This means that the
value of r does not change. The result is that the reduced mass has a circular orbit with
r = r0 . When E = E1 < 0 (or more generally for E0 < E < 0), the particle rolls back and
forth between r = r1 and r = r2 . We know from Chapter 5 that such a “rolling back and
forth” means oscillatory motion. The reduced mass gets no closer to the force center than
r = r1 , and no farther than r = r2 . Hence the resulting motion is an ellipse with rmin = r1
and rmax = r2 . Note that we are using rmin and rmax as the closest and farthest distance
between the reduced mass and the force center, respectively. In general, we can see from
Figure 9.5 that when E < 0, the motion of the reduced mass is bound to periodic orbits.
Furthermore, from Figure 9.5 we can see that when E ≥ 0, the motion is unbound. The
reduced mass comes in from infinity and gets as close to the force center as r3 , which is
not shown in Figure 9.5. However, one can find r3 by solving E = Veff (r3 ), where E is the
energy of the reduced mass. The solution can be found either analytically or numerically.
The unbound orbits are parabolic or hyperbolic, depending on the value of E.
i i
i i
i i
Table 9.1: The energies and eccentricites associated with each type of orbit for the reduced
mass µ.
Figure 9.5: A plot of the effective potential Veff (r) = VC + V (r) as a function of the distance
r. The dashed lines indicate the repulsive positive centrifugal potential VC and the attractive
potential V (r) = −2.5/r. When E = E0 the motion is a circle and r = r0 . When E = E1 < 0,
the motion is bound between the two circles with radii r = r1 and r = r2 . When E > 0, the
motion is unbound.
So far in this section, we have learned that the energy of the reduced mass determines
the nature of its orbit. Now that we have a qualitative understanding of the orbit of µ, we
can use our knowledge of conic sections in order to determine the conditions under which
each orbit occurs. We know that conic sections have specific eccentricities, for example,
e = 0 for a circle. Using (9.4.8) we can now determine the energies of each type of orbit.
They are summarized in Table 9.1.
i i
i i
i i
Figure 9.6: The motion of the planets around the Sun can be described by an ellipse which
has the Sun at one focus, F . The ellipse can be described mathematically using (9.4.5) in
polar coordinates (r, θ).
At θ = 90◦ and at θ = 270◦ the distance is equal to r = α, which is called the semi-latus
rectum. This quantity has dimensions of length, and depends on both the semi major axis
a of the elliptical motion, and on the eccentricity e:
α = a 1 − e2 (9.6.4)
i i
i i
i i
α k
a= = (9.6.5)
1−e 2 2 |E|
α `
b= √ =p (9.6.6)
1−e 2 2µ |E|
Furthermore, we can rewrite the aphelion and perihelion distances as:
Figure 9.7: Kepler’s first law places the Sun at the focus of an elliptical orbit
The next two examples shows how to plot the orbits of planets. The first example uses
data from the planets Mercury and Earth, while the second uses a numerical solution to
Newton’s second law.
Example 9.4: The orbits of Earth and Mercury around the Sun
The Earth’s distance from the Sun ranges from 147.5 million km (at perihelion), to
about 152.6 million km (at aphelion), while for Mercury these distances are from 46,000,000
to 70,000,000 km. Plot the orbits of Mercury and Earth around the Sun.
Solution:
The semi-major axis a can be found directly from the distances rmin and rmax for the
perihelion and aphelion by using (9.6.3), and the eccentricity e can be found by using
(9.6.7). Once the geometrical properties a, e are known, we know everything about the
orbits and we can plot them using the polar plot commands with the general equation
(9.6.1) for the elliptical orbits.
The Mathematica code uses the PolarPlot command to plot the two orbits and also
prints out the values of e for the two planets. In this polar plot the Sun is of course
at the origin (0, 0). Notice the almost completely circular orbit of the Earth (because the
i i
i i
i i
eccentricity e = 0.0167 is close to zero), and the clearly elliptical orbit for Mercury (because
of the larger value of e = 0.21).
a = (rmin + rmax)/2;
eEarth = (rmax − rmin)/(rmin + rmax)//N ;
a = (rmin + rmax)/2;
eMercury = (rmax − rmin)/(rmin + rmax)//N ;
OUTPUT:
The eccentricity of Earth’s orbit is: 0.0169943
The eccentricity of Mercury’s orbit is: 0.206897
Show[gr1,gr2]
2 × 1011
1 × 1011
y, m
0
Out[ ]=
-1 × 1011
-2 × 1011
-2 × 1011-1 × 1011 0 1 × 1011 2 × 1011
x, m
i i
i i
i i
d2 r dv GM m
m 2
=m = − 2 r̂
dt dt r
Following a procedure similar to Example 9.1, we obtain the following two coupled differ-
ential equations:
d2 x GM
=− x (9.6.8)
dt2 (x + y 2 )
2 3 /2
d2 y GM
=− y (9.6.9)
dt2 (x + y 2 )
2 3 /2
In the Mathematica code that follows, the command NDSolve is used to solve the system
of differential equations (9.6.8) and (9.6.9), and the solutions x(t) and y(t) are stored in
the parameters x1 and y1 correspondingly. The Table command is used to construct the
pairs of position coordinates x(t) and y(t) from t = 0 to t = 3 s, in time steps of dt = 0.01 s
with the pairs stored in the parameter xyList.
The code calculates the semi-major axes a by evaluating the maximum and minimum
distance of the planet from the origin, and also calculates the eccentricity e and the location
of the two foci of the ellipse along the x-axis. The notation [[ ]] in Mathematica means
Part. The statement xyList[[All, 1]] tells Mathematica to take the first element of every
item in the list xyList. Correspondingly, the statement xyList[[All, 2]] tells Mathematica
to take the second element of every item in the list xyList. Finally a graph of the orbit
is produced by using the ListPlot command, and the location of the foci of the ellipse is
plotted by using the Graphics and Disk commands. Note that the points along the orbit
are sparse near the perihelion (0, 0) , but they are spaced much closer near the aphelion
(1, 0). This is because the planet moves much faster at perihelion than at aphelion.
i i
i i
i i
x1 = x/.sol[[1]];
y1 = y/.sol[[1]];
xyList = Table[{x1[t], y1[t]}, {t, 0, 10, .01}];
OUTPUT:
eccentricity=0.428563
location of focus of ellipse=0.0000153831
location of focus of ellipse=0.857142
y
0.4
0.2
Out[ ]=
x
0.2 0.4 0.6 0.8 1.0
-0.2
-0.4
` = µr2 θ̇ (9.6.10)
i i
i i
i i
rdθ
dθ
r
Figure 9.8: A small wedge of the planet’s orbit traced out in a time dt.
From this magnitude we can now derive Kepler’s second law. We begin by considering a
small wedge of the planet’s orbit traced out in a time dt, as shown in Figure 9.8.
The wedge shown in Figure 9.8 forms an infinitesimal area element dA which in polar
coordinates can be written as:
r2
dA = dθ (9.6.11)
2
By dividing with dt and using (9.6.10), we obtain Kepler’s second law:
dA r2 dθ `
= = = constant (9.6.12)
dt 2 dt 2µ
Note that Kepler’s second law holds for any central force. Equation (9.6.12) is the
mathematical form of Kepler’s second law which says that a line joining a planet to the Sun
sweeps out equal areas in equal times. This is illustrated in Figure 9.9. The result of Kepler’s
second law is that the planet’s speed increases as it reaches perihelion and decreases as it
approaches aphelion.
Figure 9.9: Kepler’s second law: A line connecting a planet to the Sun (S) sweeps out equal
areas in equal times. If the time TAB required to travel distance AB is equal to the time
TCD required to travel distance CD, then the swept out area (AREA)ABS must be equal
to the corresponding area (AREA)CDS .
Example 9.6 shows how to evaluate the angular momentum of the Earth as it orbits
around the Sun, and also how to estimate the Earth’s orbital velocity.
i i
i i
i i
Example 9.6: The orbital angular momentum of the Earth around the Sun
The speed of Earth at perihelion is 30,300 m/s, and its distance from the Sun ranges
from 147.3 million km (at perihelion), to about 152.6 million km (at aphelion). Evaluate:
(a) The orbital angular momentum of the Earth around the Sun.
(b) The orbital velocity of the Earth at aphelion.
(c) The speed of the Earth at the latus rectum point α of its elliptical orbit.
Solution:
(a) When the Earth is at the aphelion (or at the perihelion) of its orbit, the velocity
vector v = dr/dt is perpendicular to the position vector r, so that the radial component
of the velocity is zero (dr/dt = 0), therefore the magnitude of the angular momentum is
found from:
dr dr
= m | r || | sin 90◦ = mrv
| ` |= m r ×
dt dt
where v represents the magnitude of the orbital velocity at these two points and m =
5.97 × 1024 kg is the mass of the Earth. By using the given value of the speed at perihelion
and the given distance from the Sun, we can find the magnitude of the angular momentum
in SI units:
vaphelion = `/ m raphelion
which gives the value of vlatusrectum = 29, 800 m/s. Note that the velocity vector at the
latus rectum is not perpendicular to the position vector.
i i
i i
i i
2µ
dt = dA (9.6.13)
`
2µ
P= πab (9.6.14)
`
where P is the period, the time needed for a line joining the reduced mass and the force
center to sweep out the entire area of the ellipse. Next, we use (9.6.5) and (9.6.6) to express
the period P and as a function of energy E.
r
µ
P = πk |E|−3/2 (9.6.15)
2
By substituting now |E| = k/(2a) from Eq.(9.6.5), we obtain Kepler’s third law:
4π 2 µ 3
P2 = a (9.6.16)
k
Kepler’s original statement was that the planet’s period squared was proportional to the
cube of its semimajor axis. Kepler didn’t know the proportionality constant. However, using
Newton’s second law, we have found it! If we consider the central force to be Newton’s law
of universal gravitation, the constant k = Gm1 m2 and the definition of the reduced mass
µ = m1 m2 / (m1 + m2 ), then Kepler’s third law becomes
4π 2 a3
P2 = (9.6.17)
G (m1 + m2 )
However, for a planet of mass m1 orbiting the Sun of mass m2 = M m1 , we have
4π 2 a3
P2 = (9.6.18)
GM
Example 9.7 illustrates Kepler’s third law with the Moons of Jupiter.
i i
i i
i i
Solution:
The data can be analyzed as shown in the Mathematica example below. First we set up
two lists called aValues and periods, which contain the semi-major axes a and the periods
P of the 4 Moons. The aValues list is converted in SI units by converting the units of
kilometers into meters, and the periods P are converted into SI units by converting days
into seconds. The command Transpose is used to create a new list called list, containing
the pairs of (a3 , P 2 ) for the four Moons. This is plotted using the ListPlot command, and
the graph is stored in the variable gr1. By using the FindFit command we fit a best line to
the graph, and we find the best slope of this line to have a slope of 3.11 × 10−16 (SI-units).
From Kepler’s law of the periods we have
4π 2 3
P2 = a
GM
where M is the mass of Jupiter. From this equation, the slope of a plot of P 2 vs a3 will
be equal to 4π 2 /(GM ).
By setting the best slope 3.11×10−16 = 4π 2 /(GM ) and using the NSolve command, we
obtain the mass of Jupiter as M = 1.89 × 1027 kg. This value is very close to the accepted
value of the mass of Jupiter MJ = 1.90 × 1027 kg.
gr1 = ListPlot[list, Frame → True, FrameLabel → {“a∧ 3 (m∧ 3)”, “P∧ 2 (s∧ 2)”},
PlotRange → All, LabelStyle → Large];
Print[“Mass of Jupiter = ”]
NSolve[slope == 4 ∗ Pi∧ 2/(G ∗ M ), M ]
OUTPUT:
Mass of Jupiter =
{{M → 1.89859 × 1027 }}
i i
i i
i i
●
2.0 × 1012
1.5 × 1012
P2 (s2)
1.0 × 1012
Out[ ]=
5.0 × 1011
●
●
0 ●
0 1 × 1027 2 × 1027 3 × 1027 4 × 1027 5 × 1027 6 × 1027
a3 (m3)
m2
r1 = r=m (9.7.1)
m1 + m2
−m1
r2 = r = 1−m (9.7.2)
m1 + m2
i i
i i
i i
y
m3
R m2
r2
t
m 1-m x
r1
m1
Figure 9.10: The coordinate system for the planar restricted circular three-body problem.
The center of mass for m1 and m2 is at the origin. Mass m1 follows a circular orbit with
radius m about the center of mass and m2 follows a circular orbit of radius 1 − m. The third
body, m3 is located at the position R. Not shown in this figure is a vector r which points
from m1 to m2 and passes through the origin.
R = X î + Y ĵ (9.7.5)
We will use the Lagrangian to derive the equations of motion for the particle m3 . The
kinetic energy of the mass m3 is,
1
T3 = m3 Ẋ 2 + Ẏ 2 (9.7.6)
2
We find the potential energy using the formula, V = −Gm1 m2 /r where m and m2 are
the masses of the interacting particles and r is the distance between m and m2 . There are
two gravitational interactions that involve m3 , V13 is the potential energy between m1 and
m3 and V23 is the potential energy between m2 and m3 . They can be written as:
G(1 − m)m3
V13 = − (9.7.7)
r13
Gmm3
V23 = − (9.7.8)
r23
i i
i i
i i
where
2
r13 = (X + m cos t)2 + (Y + m sin t)2 (9.7.9)
2 2
2
r23 = (X − (1 − m) cos t) + (Y − (1 − m) sin t) (9.7.10)
1 (1 − m)m3 mm3
L = m3 Ẋ 2 + Ẏ 2 + + (9.7.11)
2 r13 r23
Note that (9.7.11) is time-dependent. We can eliminate the time-dependence by trans-
forming our coordinate system into a frame that is rotating with the masses m1 and m2 .
In this rotating frame m1 and m2 are at rest, with m1 located at the point (m, 0) and m2
located at the point (1 − m, 0). Let us use the notation (x, y) to describe the position of m3
in the rotating frame. The transformation that connects the coordinates (X, Y ) and (x, y)
is the rotation matrix
cos t − sin t
X x
= (9.7.12)
Y sin t cos t y
where the angle of rotation is θ = ωt, with ω = 1 due to our assumption about the period
of revolution for masses m1 and m2 . We can write out the matrix multiplication in (9.7.12)
to find,
Inserting equations (9.7.15) and (9.7.18) into (9.7.11) and after some algebra, we can
rewrite the Lagrangian as,
1 m (1 − m) mm
3 3
L = m3 (ẋ − y)2 + (ẏ + x)2 + + (9.7.19)
2 r13 r23
Finally, to get the equations of motion, we use the Euler-Lagrange equations,
i i
i i
i i
d ∂L ∂L
− =0 (9.7.20)
dt ∂ ẋ ∂x
d ∂L ∂L
− =0 (9.7.21)
dt ∂ ẏ ∂y
which result in
(1 − m) (x + m) m [x − (1 − m)]
ẍ − 2ẏ − x = − 3/2 − 3/2 (9.7.22)
(x + m)2 + y 2 (x − (1 − m))2 + y 2
(1 − m) y my
ÿ + 2ẋ − y = − 3/2 − 3/2 (9.7.23)
(x + m)2 + y 2 [x − (1 − m)]2 + y 2
The equations of motion (9.7.22) and (9.7.23) are two coupled second-order ordinary
differential equations that must be solved numerically. To compute the numerical solution in
Python, we rewrote (9.7.22) and (9.7.23) as a system of four first-order differential equations
using u ≡ ẋ and v ≡ ẏ as new variables. The resulting equations of motion are
ẋ =u (9.7.24)
(1 − m) (x + m) m [x − (1 − m)]
u̇ =2v + x − 3/2 − 3/2 (9.7.25)
(x + m)2 + y 2 (x − (1 − m))2 + y 2
ẏ =v (9.7.26)
(1 − m) y my
v̇ = − 2u + y − 3/2 − 3/2 (9.7.27)
2
(x + m) + y 2 [x − (1 − m)]2 + y 2
The four first-order equations of motion are then solved using Algorithm 8. In Algorithm
8, odeint is used to solve equations (9.7.24) to (9.7.27) for the condition m = 0.01. For
more information on odeint, see Section 2 in Chapter 13. The local variables peXterm and
peYterm are the right hand side of (9.7.22) and (9.7.23), respectively. These additional
variables were included to improve the clarity of the algorithm. In addition, we included the
basic information needed to create a four-panel plot. The results of Algorithm 8 are shown
in Figure 9.11. Note that to create Figure 9.11, you need to include additional formatting to
include the dots and the plot titles. Further note that we used line wrapping in Algorithm
8 to prevent the code from overrunning the margins. The hook arrows represent the line
break, you will not include them in your own code.
In Figure 9.11, we see the trajectory of m3 in the rotating coordinate frame (x, y) for
four different initial conditions. In each plot, the red (in the e-book) dot near the origin is
m1 and the other dot (which is green in the e-book) is m2 . The blue (in the e-book) curve
is the trajectory of m3 . Notice in the upper-left plot of Figure 9.11, m3 is in a precessing
orbit about m1 , which is obscured by the trajectory. The upper-right plot has m3 spiraling
out towards an orbit that encompasses both m1 and m2 . The lower-right plot shows a case
where m3 orbits about m2 .
Arguably the most interesting plot in Figure 9.11 is the lower-left plot, which shows
an example where m3 orbits around a so-called Lagrange point of the m1 -m2 system. A
Lagrange point is a point in space where the centrifugal force is balanced by the gravitational
i i
i i
i i
import numpy a s np
from s c i p y . i n t e g r a t e import o d e i n t
import m a t p l o t l i b . p y p l o t a s p l t
mu = 0 . 0 1
d e f threeBody ( vec , t ) :
x , y , u , v = vec
i n i t i a l C o n d i t i o n 1 = [ −0.75 , 0 , 0 , 1]
t=np . l i n s p a c e ( 0 , 2 0 , 2 0 0 1 )
s o l 1 = o d e i n t ( threeBody , i n i t i a l C o n d i t i o n 1 , t )
ax1 . p l o t ( s o l 1 [ : , 0 ] , s o l 1 [ : , 1 ] )
ax2 . p l o t ( s o l 2 [ : , 0 ] , s o l 2 [ : , 1 ] )
ax3 . p l o t ( s o l 3 [ : , 0 ] , s o l 3 [ : , 1 ] )
ax4 . p l o t ( s o l 4 [ : , 0 ] , s o l 4 [ : , 1 ] )
p l t . show ( )
i i
i i
i i
Figure 9.11: The results of Algorithm 8. The red (in the e-book) dot near the origin is m1
and the green (in the e-book) dot is m2 . The curve represents the path followed by m3 .
attraction. Because we transformed the coordinates of the original problem into a rotating
system, we introduced a non-inertial frame. We will discuss non-inertial frames in the next
chapter. However, for now all you need to know is that in non-inertial frames, an apparent
force called the centrifugal force appears which is directed outward. Near two orbiting
masses, such as m1 and m2 , there are five points called Lagrange points where the centrifugal
(outward) force is equal to the gravitational attraction (inward). The Lagrange points are
equilibrium points and the one represented in the lower left panel of Figure 9.11 is called L4,
which is a stable equilibrium. There are asteroids called Trojan asteroids that orbit Jupiter’s
L4 point. In addition, Lagrange points in Earth-Sun and Earth-Moon systems have been
used by many probes and satellites. For example, the James Webb Space Telescope is
planned to orbit L2, which lies along a line joining the Earth and the Sun, and would be to
the right of the green dot in Figure 9.11.
F = f (r)r̂
where the vector r is the location of the particle relative to the force center. Central forces
conserve energy and angular momentum.
In a two-body system, where two masses m1 and m2 interact via a central force, we
can reduce the problem to that of the motion of a single particle of mass
i i
i i
i i
m1 m2
µ=
m1 + m2
called the reduced mass. The reduced mass is located at r = r1 − r2 relative to the system’s
center of mass. The center of mass serves as the force center.
Although there are many different approaches to finding the equations of motion for
the reduced mass in the two-body problem, one of the easiest is by solving the differential
equation,
d2 1 1 r2
2
+ = − 2 f (r)
dθ r r µ`
From the equation of motion, which is derived from Newton’s second law in polar
coordinates, we can derive general forms of Kepler’s three laws of planetary motion.
The first law is that the orbit of the reduced mass is of the form of a conic section with
the force center at one focus,
s
1 1 `2 2E`2
= (1 + e cos θ) α= e = 1+
r α µk µk 2
For planets orbiting the Sun, their motion is an ellipse with the Sun at one focus. The
shape of the ellipse can be described by
rmin = a (1 − e) rmax = a (1 + e)
where e is the eccentricity of the ellipse, r is the distance from the Sun to the planet, and
a and b are the semi-major and semi-minor axes of the ellipse, respectively. The perihelion
is the point where the distance is minimum rmin , and the aphelion is the point where
the distance is maximum rmax . In general, the motion of the reduced mass can be any
conic section. The choice of conic section is determined by the energy which determines
the eccentricity e of the orbit. If 0 < e < 1 the motion is an ellipse, if e = 0 the orbit is a
circle, if e > 1 the motion is a hyperbola, and for e = 1 it is a parabola.
The conservation of angular momentum leads to Kepler’s second law, which states that
a line joining the reduced mass to the force center sweeps out equal areas in equal time
periods.
Kepler’s third law is obtained by integrating the mathematical form of Kepler’s second
law and it relates the period of the planet’s orbit P to the semimajor axis of its orbit a.
4π 2 a3
P2 =
G (m1 + m2 )
Finally, the three-body problem involves interactions between three objects. We exam-
ined the circular restricted three-body problem which is a relatively simple formulation of
the three-body problem. However, even the circular restricted three-body problem must
be analyzed numerically. The motion of the third mass m3 m1 , m2 can be rather com-
plicated.
i i
i i
i i
4. A particle moves in a central force field defined by F = −Kr2 . It starts from rest at
a point on the circle r = a.
(a) Prove that when it reaches the circle r = b its speed will be
q
v = 2K(a3 − b3 )/3m
5. A particle of mass m moves in a central force field F = −K/rn where K and n are
constants and K > 0. It starts from rest at r = a and arrives at r = 0 with finite speed
Vo .
(a) Prove that we must have n < 1.
(b) Prove that s
2Ka1−n
Vo =
m (n − 1)
6. Let L, M, and T represent the dimensions of length, mass and time, respectively. Find
the dimensions of the universal gravitational constant G.
i i
i i
i i
7. The coordinates of a mass m are given as a function of time t, as x = x0 cos (ω1 t) and
y = y0 sin (ω2 t).
(a) Show that the force is a central force if ω1 = ω2 .
(b) Show that the total energy is conserved.
8. A mass m experiences a force F = f (r)r̂ − λv, where λ = positive constant, f (r) is
a scalar function of the radial distance r, v is the velocity vector, and r̂ is the unit
vector along the radial direction. Show that the angular momentum varies with time
λ
as ` = `0 e− m t , where `0 is the initial angular momentum at t = 0.
Section 9.2–9.3: The Two-Body Problem and Equations of Motion for the Two-
Body Problem
9. Two point masses m and M are separated by a distance r. Starting from rest, they
both begin to accelerate towards each other. Show that the elapsed time t until they
collide is: s
r3
t=π
8G(M + m)
10. Show that the total angular momentum L =``1 + ` 2 of the two mass system is equal to
L =µ (r × v), where µ is the reduced mass, r = r2 − r1 is the relative position vector
of the two masses, and v = ṙ.
11. Show that the kinetic energy of the two mass system with a reduced mass µ and
relative position vector of the two masses r = r2 − r1 is equal to:
1
T = µ | ṙ |2
2
12. In Example 9.1 we evaluated numerically and plotted the orbits in a two-body gravi-
tational system. Use the command ListAnimate in Mathematica, to produce an ani-
mation of the two bodies moving around each other.
13. A particle moving in a central force whose center is located at r = 0, describes a spiral
r = e−θ . Prove that the force is proportional to 1/r3 .
14. What central force is required to produce an orbit r2 = a2 cos (2θ)?
15. Show that both r = e−θ and r = 1/θ are possible orbits for a central force proportional
to 1/r3 . How is this possible?
16. Show that an inverse cube central force is required to make a particle move around
the origin O with a speed inversely proportional to the distance from O. What types
of orbits are possible?
17. The orbit of a mass m inside a central force field is such that r = a/θ where a =
constant. Show that the potential must be proportional to 1/r2 .
18. Obtain the orbit for a mass m moving in a central force field defined by F = −K/r3 ,
and describe it physically.
i i
i i
i i
19. A particle of mass m moves in a central force field given in magnitude by F = −Kr
where K is a positive constant and r is the position vector. If the particle starts at
r = a, θ = 0 with a speed v0 in a direction perpendicular to the x-axis, determine and
describe its orbit.
20. A particle of mass m moves in a central force field given in magnitude by F = −K/r3
where K is a positive constant, and r is the distance from the force center. The particle
starts at r = a, θ = 0 with a speed v0 in a direction making an angle α with the positive
x axis.
(a) Show that the differential equation for the orbit is given in terms of u = 1/r by
d2 u
+ (1 − γ)u = 0
dθ2
where
K
γ=
ma2 v 20 sin2 α
(b) Solve the differential equation in (a) and interpret it physically.
21. Show that if the force is given by
A B
F= or F = − 2 3
r4 cos θ r cos θ
then a possible orbit in both cases is given by r = 2a cos θ.
(a) What can we conclude about uniqueness of forces when orbits are specified?
(b) Answer part (a) when the forces are central forces.
y 2 x2
− =1
b2 a2
Graph this equation with a = b = 1, find the vertices, the foci and the directrices.
i i
i i
i i
26. The distance of closest approach of Halley’s comet to the Sun is 0.57 astronomical
units AU (1 AU is the mean Earth-Sun distance). The greatest distance of the comet
from the Sun is 35 AU. Do a polar plot of the orbit of Halley’s comet around the Sun,
in the vicinity of our solar system, which is about 40 AU. Plot together the orbits of
Mercury, the Earth and Halley’s comet.
Section 9.5: Orbits in Central Force Fields
27. Discuss the motion of a particle in the central force field F = α/r2 + β/r3 for β > 0
28. If a particle moves in a circular orbit under the influence of a central force whose
center is at the center of the circular orbit, prove that its speed around the orbit must
be constant.
29. Planet A has a velocity vA at aphelion and a velocity of vP at perihelion. A second
planet B moves around the same star in a circular orbit with a radius R and with a
speed v. Show that the aphelion distance for planet A is
2Rv 2
RA =
vA (vA + vP )
30. A comet of energy E and angular momentum ` enters a region with an attractive
central potential V (r). Show that the distance of closest approach is
`
rmin = p
2m (E − V (rmin ))
31. A comet of mass m approaches our Sun with a velocity v0 from a very large distance
as shown in Figure 9.12. If there was no deflection of the comet, it would have passed
a distance b from the Sun as shown. Show that the distance of closest approach for
the comet is given by: s 2
k k
rmin = + + b2
mv02 mv02
where the gravitational potential is V = −k/r.
Figure 9.12: Problem 9.31: The comet of mass m approaching the Sun.
i i
i i
i i
32. A spacecraft travels from planet A to planet B in an elliptical orbit, so that the
perihelion is at planet A and the aphelion is at planet B. Planets A and B move in
circular orbits around a star of mass M , with radii RA and RB as shown in Figure
9.13. Assume that the masses of the planets have no effect on the spacecraft. Show
that the spacecraft must be given a speed
s
2GM RA
r
GM
v= −
RB (RA + RB ) RA
with respect to planet A in order to move in the dashed elliptical orbit shown below.
Figure 9.13: Problem 9.32: A spacecraft transiting from planet A to planet B along the
dotted elliptical path.
i i
i i
i i
39. Find the distance of the Moon from the Earth, by assuming a circular orbit around
the stationary Earth.
40. Calculate the mass of the Sun using the fact that the Earth is 150 million km from
it, and it makes one complete revolution about it in approximately 1 year.
41. In this exercise you will study Kepler’s third law using the databases available in
Mathematica.
(a) Type the PlanetData[] command, this will load the data about the planets from
the Mathematica database. It may take about 1 minute or so to load all the data.
Next define a list with the planets, using a variable called planetsList, so we can use
it further down in our calculations
planetsList = PlanetData[]
(b) Now type the command: PlanetData[planetsList, Properties]
This shows you all the different properties of the planets for which there is available
data. Hover with your mouse over the name of the different properties, for example
by hovering over orbital period you can see underneath OrbitPeriod for the name of
this property.
Now get the orbital period of all the planets: PlanetData[planetsList, OrbitPe-
riod]
This gives us the orbital period of all the planets in years. This is not an SI unit, so
let us change it into SI units of seconds (s) by using the UnitConvert command,
UnitConvert[PlanetData[planetsList ,OrbitPeriod]]
This looks better in SI units, so let’s change these periods to SI units, and also call
this a new list periodsSI :
periodsSI = UnitConvert[PlanetData[planetsList ,OrbitPeriod]]
(c) Create a list called semimajorSI, which will have the semi-major axis of all the
planets’ orbits around the Sun, in SI units of course!
(d) Now that you have two lists containing the periods and the semi-major axes,
create pairs of these using the Transpose command. Finally plot this new list using
the ListPlot command. It looks like this graph is not a straight line, because it does not
represent Kepler’s Third Law. Change your graph so that it does represent Kepler’s
Third Law, and it is also a straight line!
42. Using the code provided in Section 9.7, try different initial conditions for the mass
m3 and discuss the result.
43. In this problem, we will plot the Lagrange points for the circular restricted three-
body problem with m = 0.01. We continue with the situation where m1 m2 such
as the Earth-Sun-satellite system. Note that the value of m in the Earth-Sun satellite
is close to 10−6 but we will continue with m = 0.01 for visualization purposes. The
Lagrange points are points where ẍ = ÿ = 0 and ẋ = ẏ = 0 in the rotating frame
i i
i i
i i
discussed in Section 9.7. Using a computer, solve numerically (9.7.22) and (9.7.23)
when ẍ = ÿ = ẋ = ẏ = 0. Plot the resulting points on the same graph as a circle with
radius 1 − m, the orbit of m2 . You should find that there are a total of five Lagrange
points and that three of the Lagrange points lie along the orbit of m2 and three of
them lie along the same line (one of those are also on the orbit).
i i
i i
i i
CHAPTER 10
Motion in Noninertial
Reference Frames
As we saw in Chapter 1, Newton’s second law is valid only for inertial reference frames.
However, it is sometimes the case where we need to describe the motion in a noninertial
reference frame. For example, long-range missile trajectories need to be described using
a noninertial reference frame, a frame fixed to the surface of the rotating (and therefore
noninertial) Earth. In Chapter 1 we found that using Newton’s second law in a noninertial
reference frame, resulted in the appearance of inertial forces . In this chapter, we will expand
on the idea of inertial forces; specifically, those arising in rotating reference frames. We will
discuss how a vector in a rotating reference frame can be described in a nonrotating frame,
and from that description, show how applying Newton’s second law in a rotating frame
results in inertial forces such as the Coriolis force and the centrifugal force. We will then
study how the Coriolis and centrifugal forces affect the motion of a particle near the surface
of the Earth. Finally, we will explore the famous problem of the Foucault pendulum and
projectile motion in a noninertial frame.
ṙ0 = ṙ + V (10.1.1)
where V = Ṙ is the instantaneous velocity of the frame S relative to S0. For example, if
S 0 is a stationary reference frame on the ground and S is fixed in an accelerating car, then
the particle’s velocity relative to the ground (ṙ0 ) is the velocity of the particle relative to
311
i i
i i
i i
y m
r
y' r'
S x
R
S' x'
Figure 10.1: A particle of mass m located at a position r0 relative to the origin of an inertial
frame S 0 and at a location r relative to the origin of a noninertial frame S. The vector R
measures the position of the noninertial frame’s origin relative to the origin of the inertial
frame.
the car (ṙ) plus the velocity of the car (V). It should be noted that (10.1.1) only holds for
non-relativistic cases.
Next, we differentiate (10.1.1) to get the acceleration of the particle as measured in the
noninertial frame,
r̈ = r̈0 − A (10.1.2)
where A = V̇. Notice that similar to what was found in Chapter 1, an additional acceleration
measured in S 0 that is not measured in S. Multiplying (10.1.2) by the mass of the particle
m,
where F = mr̈ are the forces acting on the particle as measured in the noninertial frame,
and F0 = mr̈0 are the forces acting on the system as measured by the noninertial frame. The
forces acting on the particle may include gravity, air resistance, friction, etc. These forces are
measured in both frames. However, there is an additional force measured in the noninertial
frame, −mA. This additional force is sometimes called an inertial force, a fictitious force,
or a pseudo force. In order to use Newton’s second law to describe motion in a noninertial
frame, we need to include the additional inertial force, which is not due to the particle’s
interaction with another body, but rather arises from the acceleration of the noninertial
frame.
Although the inertial force is sometimes called the fictitious force, inertial forces are real
to those moving in a noninertial reference frame. You have experienced noninertial forces
in an accelerating car as you are being pushed back on the seat, or on a rotating carnival
ride such as the Tilt-a-Whirl, where you are pushed against the side of a carriage as the
ride spins. In each case, the force is due to your inertia as opposed to the interaction with
another body, hence the name inertial since it is introduced in order to account for all of the
accelerations you are experiencing. By including inertial forces, we are able to use Newton’s
second law to describe motion in noninertial reference frames.
i i
i i
i i
θ ℓ
Solution:
According to (10.1.3), the force in the frame S which is accelerating with the pendulum
is,
F = (T + mg) − ma (10.1.5)
where T is the tension in the pendulum, g is the acceleration due to gravity, and a is a
vector that has a magnitude of a and points in the direction opposite of g. The term in
the parentheses of (10.1.5) is F0 , the sum of the forces acting on the system according to
Newton’s second law. We can rewrite (10.1.5) as
F = T + mgeff (10.1.6)
where geff = g −a is the effective gravitational field which includes the gravitational attrac-
tion and the acceleration of the pendulum’s support. It then follows that in the frame S,
the equation of motion of the pendulum for small oscillations is,
geff
θ̈ + θ=0 (10.1.7)
`
where θ is the angle that the pendulum makes with the vertical (as usual). From Chapter
6 we know that the coefficient in front of θ in (10.1.7) is the square of ω0 , the frequency
of small oscillations. Therefore, we find
r
|g − a|
ω0 = (10.1.8)
`
We see that the upward acceleration effectively reduces the gravitational field. The
result is that the accelerated frame is the same as having an additional gravitational
field (in this case, whose acceleration is in the opposite direction of g). The fact that
accelerating frames are indistinguishable from gravitational fields is an important element
in the general theory of relativity.
i i
i i
i i
Figure 10.2: The right-hand rule. The circle represents a wheel and the arrows on the circle
represent a counter clockwise rotation of the wheel. Curl the fingers of your right hand in
the direction of the arrows on the circle. Your thumb should point in the direction of the
vector ω .
infinitesimal displacement of a rigid object such that a point on the body remains fixed, is
equivalent to a rotation about an axis that runs through the fixed point. This theorem is
difficult to prove, but we don’t need to prove it here. However, as an illustrative example,
consider a wheel rolling along a road. Euler’s theorem says that an infinitesimal displacement
of the wheel can be described as a rotation about the contact point between the wheel and
the road.
Euler’s theorem tells us that in order to specify the rotation about a point, we need
only the direction of the axis and the amount rotated about the axis. Of course, if we are
interested in the rate of rotation, the angular velocity, then we’d need the direction of the
axis of rotation and the rate of the rotation. This means we can write the angular velocity
as a vector ω that lies along the axis of rotation and whose magnitude is the rate of rotation.
A stationary spinning top may have an angular velocity of 2π rad/sec with an axis that is
vertically oriented. However, does ω point up or down? The answer to that is determined
by the right-hand rule,
Curl the fingers of your right hand in the direction of the rotation. Your thumb points in
the direction of ω .
The right-hand rule is illustrated in Figure 10.2. The circle with arrows in Figure 10.2
represents a wheel rotating in a counter clockwise direction. As you curl your fingers in the
direction of the arrows on the wheel, your thumb should point up, away from the page (or
screen if you are reading this on a device). Your thumb is pointing in the direction of the
wheel’s angular velocity.
If an object’s angular velocity is changing with time, then both the rate of rotation
and/or the orientation of the rotation axis is changing in time. As you might imagine,
situations that involve changes in the orientation of the rotation axis can be difficult to
describe mathematically. In this chapter, we will focus primarily on systems with fixed
angular velocities, i.e., both the rate of rotation and the direction of the rotation axis will
remain constant.
i i
i i
i i
ρ v
θ r
Figure 10.3: The northern hemisphere of the Earth showing the location of a point particle
(black dot) and its velocity, v.
v = ω ×r (10.2.1)
Notice that (10.2.1) is the general form for the equation v = rω you learned in introductory
physics (where θ = π/2). We need to use (10.2.1) in this case because the origin of the
inertial frame is the center of the Earth, and the center of the Earth is not the center of
the circle along which the particle travels. Furthermore, the center of the Earth will serve
as the origin of the inertial frame that we will use in future problems.
Equation (10.2.1) is not unique to the velocity vector. Note that we can rewrite (10.2.1)
as,
dr
= ω ×r (10.2.2)
dt
which is not a formula unique to the vector r. In fact, for any vector Q that is constant in
the rotating reference frame, its time derivative as measured in the non-rotating frame is,
i i
i i
i i
ê3
ω P
ê'3
r ê2
r'
R ê1
ê'2
ê'1
Figure 10.4: The inertial frame S 0 (primed unit vectors) and the noninertial rotating frame
S (unprimed unit vectors, red in e-book). The vector R represents the location of the origin
of the noninertial frame as measured in the inertial system. The vector ω represents the
angular velocity of the noninertial frame.
dQ
= ω ×Q (10.2.3)
dt
In the next section, we will provide a derivation of a more general form of (10.2.3) which
will include vectors that are not constant in the rotating reference frame.
i i
i i
i i
Now let’s consider an arbitrary vector Q which is measured in the noninertial frame.
The vector Q could be the position of a particle at the point P , the velocity of the particle,
or any other vector quantity. We will compute the time derivative of Q as measured in the
inertial (fixed) frame and the noninertial (rotating) frame using the following notation,
dQ
= the time derivative of Q relative to the inertial (fixed) frame S’
dt f
dQ
= the time derivative of Q relative to the noninertial (rotating) frame S.
dt r
We begin by writing Q in terms of the unit vectors that are fixed in the noninertial
frame,
X3
Q= Qi êi (10.3.1)
i=1
The coefficients Qi are the same in each frame. Observers in the inertial frame, would
see êi vary with time, but that would be the only difference between what is observed in
the two frames. If we calculate the time derivative of Q relative to the noninertial frame,
we would find,
3
dQ X
= Q̇i êi (10.3.2)
dt r
i=1
Because the scalar coefficients Qi are the same in each frame, we do not need to distin-
guish whether Q̇i is relative to the inertial or noninertial frame. Note that the derivative
doesn’t affect the unit vectors because they are constant relative to the noninertial frame.
Next, we compute the time derivative of Q relative to the fixed frame,
3 3
dQ X X dêi
= Q̇i êi + Qi (10.3.3)
dt f dt f
i=1 i=1
The first term on the right-hand side of (10.3.3) is what we found when computing the
time derivative relative to the noninertial frame. Notice that the second term in the right-
hand side of (10.3.3) is the time derivative of the vector êi , which is constant in the rotating
frame. We can use (10.2.3) to write,
dêi
= ω × êi (10.3.4)
dt f
Therefore, (10.3.3) becomes,
3 3
dQ X X
= Q̇i êi + ω × (Qi êi ) (10.3.5)
dt f i=1 i=1
or using (10.3.2):
i i
i i
i i
Notice that (10.3.6) is what we would expect from our studies of translationally moving
reference frames. To better understand this, let Q = r, the position of a particle as measured
in the noninertial (rotating) frame. The first term in the right-hand side of (10.3.6) is the
velocity measured in the rotating frame. The second term is the translational velocity of
the particle at r due to the rotation of the noninertial frame. Hence, the velocity measured
by an observer in the fixed frame is the velocity of the particle in the moving frame plus the
velocity of the frame itself. Note that one result of (10.3.6) is that the angular acceleration
ω is the same in each frame because ω × ω = 0.
ω̇
Equation (10.3.6) tells us that a vector’s rate of change can be related between inertial
and noninertial frames. Newton’s second law provides a particle’s equations of motion by
stating that the net force acting on a particle is proportional to the particle’s change in
velocity as measured in an inertial frame. Thus, (10.3.6) allows us to relate the acceleration
of a particle in an inertial frame, found from Newton’s second law, to situations where the
observer is in a noninertial frame. In the next section, we will use (10.3.6) to find Newton’s
second law in a rotating frame.
where v0 is the velocity of the particle relative to the inertial frame S 0 and F0 is the net
force acting on the particle. The net force F0 includes the interaction forces between the
mass m and other bodies such as gravity, air resistance, friction, etc. In order to find the
form of Newton’s second law for a noninertial frame, we will need to rewrite the acceleration
as measured in the inertial frame (dv0 /dt)f in terms of the acceleration as measured in the
noninertial frame.
We begin by finding v0 . Figure 10.4 shows that,
r0 = R + r (10.4.2)
The velocity v0 is found by taking the time derivative of r0 relative to to the inertial frame.
0
0 dr dR dr
v = = + (10.4.3)
dt f dt f dt f
We use (10.3.6) on the second term in the right-hand side of (10.4.3) in order to find
how an observer in the inertial frame S 0 measures the velocity of the particle moving in the
noninertial frame S,
dr dr
= +ω ×r (10.4.4)
dt f dt r
i i
i i
i i
where ω is the angular velocity of S. By using (10.4.4) and defining the following:
0
dr
v0 ≡ (10.4.5)
dt f
dR
V≡ (10.4.6)
dt f
dr
v≡ (10.4.7)
dt r
we can write:
v0 = V + v + (ω
ω × r) (10.4.8)
where:
Now, we will go through the right-hand side of (10.4.9) term by term. The first term in
the right hand side of (10.4.9) is the linear acceleration of the noninertial frame,
dV
A= (10.4.10)
dt f
The second term in the right-hand side of (10.4.9) is rewritten using (10.3.6):
dv dv
= + (ωω × v) (10.4.11)
dt f dt r
=a + (ω
ω × v) (10.4.12)
where a = (dv/dt)r is the acceleration of the particle as measured in the noninertial frame.
The third term in the right hand side of (10.4.9) does not need simplification because ω̇ is
the same in each reference frame. Finally, the fourth term in the right hand side of (10.4.9)
is rewritten using (10.3.6),
dr dr
ω× =ωω× + ω × (ω
ω × r) (10.4.13)
dt f dt r
= (ω
ω × v) + [ω
ω × (ω
ω × r)] (10.4.14)
F0 = mA + ma + 2m (ω
ω × v) + m (ω̇
ω × r) + m [ω
ω × (ω
ω × r)] (10.4.15)
i i
i i
i i
If we define
F ≡ F0 − mA − 2m (ω
ω × v) − m (ω̇
ω × r) − m [ω
ω × (ω
ω × r)] (10.4.16)
to be the force experienced by the particle in the noninertial frame, then for noninertial
frames, Newton’s second law takes the form
F =ma (10.4.17)
Recall that a is the acceleration of the particle as measured in the noninertial frame.
Although the above equation looks like how we have written Newton’s second law from
Chapter 1 and onwards, it is important to remember all of the terms that are packed into
F as defined in (10.4.16). The force F is sometimes called the effective force acting on the
particle because it includes both interaction forces and inertial forces.
Notice that there are four inertial forces introduced in (10.4.16). The inertial force −mA
is the same one we found in Section 10.1 and due to the linear acceleration of S. The second
term in (10.4.16), −2m (ω ω × v), is the so-called Coriolis force which we will discuss in detail
later. The term, −m (ω̇ ω × r), is the inertial force associated with the angular acceleration of
S 0 . For most of the problems in this chapter, the frame S will be attached to the surface
of the Earth, and to a good approximation, ω̇ ω = 0. This force will not be present. The final
term, −m [ω × (ω × r)], is the centrifugal force which we will also discuss later in detail.
ω ω
We have learned that Newton’s second law can be used, with modification, for problems
involving noninertial frames. The trick is that we have to include additional inertial forces
and not just the forces of interactions between bodies in the system. Just like in Section
10.1, these inertial forces are real to the observer in the noninertial frames.
To conclude this section, we will discuss further the Coriolis and centrifugal forces. These
two intertial forces play an important role when describing the motion of an object near
the Earth’s surface.
For the remainder of this chapter, the inertial frame S 0 has its origin at the center of
the Earth and ê03 points along the Earth’s rotation axis from the South Pole to North Pole.
Hence ω = ωê03 where ω = 2π rad/day. Furthermore, the noninertial frame S has its origin
fixed to one point on the Earth’s surface, the ê1 ê2 -plane is tangent to the Earth’s surface at
the origin of S, and ê3 , which isis pointing locally up (away from the center of the Earth).
i i
i i
i i
ω ω×r
ω×r
ω ω×(ω×r) -ω×(ω×r)
ω×(ω×
ρ r)
-ω×(ω
ρ
×r)
r
θ
(a) The perspective from a distant observer. (b) A “bird’s-eye view” of Figure 10.5(a) looking
down on the vector ω .
Figure 10.5: Diagram of the centrifugal force in two perspectives. A particle is located at r
in a rotating reference frame which has an angular velocity ω. An observer in a noninertial
frame would see the particle as moving in a circle with a radius ρ = r sin θ. The particle
experiences a force −m [ω ω × (ω
ω × r)] that points radially outward from the particle’s center
of revolution.
i i
i i
i i
ω
-ω×(ω×r)
g0
ρ g
θ R
E
The position vector RE has a magnitude of 6, 371, 000 m and points in a direction from
the center of the Earth to the location of the particle on the Earth’s surface. The angle θ
is the same angle from Figure 10.5a and is related to the particle’s latitude λ such that
λ = π/2 − θ. To a stationary observer above the Earth, the particle appears to move in a
circle of radius ρ = RE sin θ.
According to (10.4.16), the effective force acting on the particle is,
F = F0 + mg0 − m [ω
ω × (ω
ω × RE )] (10.4.20)
where
GME
g0 = − 2 R̂E
RE
is the acceleration due to the gravitational attraction, found from Newton’s universal law of
gravitation, the Earth’s mass is ME = 5.972 × 1024 kg and F0 is the sum of all of the forces
other than gravity acting on the particle. The only inertial force acting on the particle is
the centripetal force. The Coriolis force is zero because the particle is not moving relative
to the Earth’s surface and, to good approximation, ω is a constant and therefore ω̇ ω = 0.
i i
i i
i i
We can collect the second and third terms in (10.4.20) and define a new acceleration
due to gravity called an effective acceleration due to gravity g where,
g = g0 − ω × (ω
ω × r) (10.4.21)
and it is g that determines the acceleration due to free fall and the period of a pendulum at
a point on Earth. Note that the centrifugal acceleration effectively reduces the acceleration
due to gravity. In addition, a plumb line will lie along the direction of g and will slightly
deviate from true vertical because of the outward-directed centrifugal force shown in the
diagram above. Furthermore, the surfaces of the Earth’s oceans are perpendicular to g not
g0 .
We can compute the effective acceleration due to gravity at the North Pole and equator
by first finding the magnitude of g0 ,
Next, we compute the magnitude of the centripetal acceleration at the North Pole and
equator. We will use acent = RE ω 2 sin θρ̂ρ which comes from dividing (10.4.19) by the mass
m. Note that we want the component acent sin θ of acent which is parallel to RE . Note
that for the North Pole, θ = 0 and for the equator θ = π/2. In each case ω = 2π rad/day =
7.272 × 10−5 rad/s,
F = F0 + mg − 2m (ω
ω × v) (10.4.22)
where F0 is the sum of all of the interaction forces experienced by the particle excluding
its gravitational attraction to the Earth. We then combined mg0 with the centrifugal force
into a new force mg. The final term in (10.4.22) is the Coriolis force, which will be the topic
for the next section.
Fcor = −2m (ω
ω × v) (10.4.23)
It is clear from (10.4.23) that the Coriolis force depends on the velocity of a particle.
Particles that are not moving will not experience a Coriolis force. We can estimate the size
i i
i i
i i
of the Coriolis force using (10.4.23). In the previous subsection, we used ω ≈ 7.3 × 10−5
rad/sec. Therefore a golf ball at a latitude of 45 degrees North traveling at 70 m/s will
experience a Coriolis force of 0.007 N, a very small effect considering its time of flight.
However, an intercontinental ballistic missile traveling at 7 km/s will experience a Coriolis
force of 0.7 N which could have a more considerable effect on the missile. However, missiles
are also affected by the weather and their courses are corrected by computer guidance
systems.
The direction of the Coriolis acceleration (and therefore Coriolis force) is determined
by the ω × v term. As shown in Figure 10.6, the Coriolis acceleration is to the right of
the velocity when the particle is moving in Earth’s Northern Hemisphere. In the Southern
Hemisphere, the Coriolis acceleration will be to the left of the particle’s velocity.
Besides missile trajectories, the Coriolis force has an important impact on Earth’s
weather, deflecting air currents in the atmosphere. In the Northern Hemisphere, air is
deflected toward the right, whichleads to high pressure to the right of the airflow, and
low pressure to the left (opposite in Southern Hemisphere), resulting in the counterclock-
wise circulation of air in the Northern Hemisphere and the clockwise circulation in the
Southern. The result can be hurricanes and typhoons.
-ω×v
Figure 10.6: The potential direction of the Coriolis acceleration, acor = −2 (ω ω × v), can
exert on a particle moving with a velocity v relative to the angular velocity ω of a rotating
reference frame. The acceleration deflects the particle towards its right in the Northern
Hemisphere and towards its left in the Southern Hemisphere.
i i
i i
i i
ez
RE
θ
λ ex
The diagram above depicts the Earth (as the circle). Fixed on the Earth’s surface is the
origin of the rotating frame S. The origin of S is located at a latitude λ. We have chosen
our coordinates in the following way. The êz -direction is locally upward, opposing the
direction of the effective acceleration due to gravity. The direction, êx , points towards the
South. Not shown is the direction êy that points towards the East in order to preserve a
right handed coordinate system. Note that the angular velocity vector ω = ωẑ0 , where ẑ0
is a unit vector in the inertial frame (fixed to the center of the Earth), points from the
Earth’s center to the North Pole.
In order to find the horizontal displacement of the falling particle, we need to solve
Newton’s second law for rotation using our established coordinate system. Newton’s second
law is,
F = mr̈ = mg − 2m (ω
ω × v) (10.4.24)
Hence, we will need to integrate Newton’s second law two times in order to get the
displacement we are seeking. Note that the vectors g and v are already in the unprimed,
noninertial, coordinate system. However, ω is in the primed, or inertial, coordinate system.
We can transform ω to the noninertial coordinate system, by thinking of the unprimed
coordinates as a rotation of the primed coordinates through an angle θ about the êy
direction. In order to obtain the form of ω in the noninertial coordinates, we can apply
the rotation matrix about the y-axis to the vector ω in the inertial coordinates,
If we assume that the x and y components of the falling particle’s velocity are negligible
and use vz = −gt, we can compute the Coriolis acceleration ω × v using the result of
(10.4.25),
êx êy êz
ω × v = −ω cos λ 0 ω sin λ = −ωgt cos λêy (10.4.26)
0 0 −gt
i i
i i
i i
Using g = −gêz and (10.4.26) we find that (10.4.24) becomes three equations,
ẍ =0 (10.4.27)
ÿ =2ωgt cos λ (10.4.28)
z̈ = − g (10.4.29)
which show us that there is no deflection to the south (x-direction) and the vertical motion
is a free fall. The deflection is towards the east (y-direction) and is found by integrating
(10.4.28) twice, once with with initial condition ẏ(0) = 0 and then a second time with
initial condition y(0) = 0. The result is,
1
y(t) = ωgt3 cos λ (10.4.30)
3
We can find the time of flight t by integrating (10.4.29) twice to get z(t) = z(0) −
1 2
(note ż(0) = 0) with z(0) = H. Solving for t, we find t = 2H/g and therefore the
p
2 gt
displacement from the plumb line of the falling particle is:
s
1 8H 3
y= ω cos λ (10.4.31)
3 g
How large is this? An object dropped from a height of 100 m at a latitude of 45◦ North
is deflected by 1.55 cm, if we neglect the effects of air resistance.
i i
i i
i i
ez
T
ey
Tx
Ty
ex
mg
Figure 10.7: A Foucault pendulum of mass m and length ` suspended from a large vertical
height at a point on the z-axis.
T
a= + g − 2 (ω × v) (10.5.1)
m
To get x(t) and y(t) (note that with our assumptions z(t) = 0), we will need to find the
components of the vector equation (10.5.1), integrate and solve the resulting second order
ODE.
To get the ODE, we write the components of each of the vectors on the right-hand side of
(10.5.1) noting that a = ẍêx + ÿêy + z̈êz . We can see in Figure (10.7), that g = −gêz where
g is the magnitude of the local effective acceleration due to gravity. Next, the components
of the tension can be found using the fact that α is a small angle and ` is large,
x y
Tx ≈ − T Ty ≈ − T Tz ≈T (10.5.2)
` `
The Coriolis acceleration ω × v can be found using the method outlined in Example
10.3,
êx êy êz
ω × v = −ω cos λ 0 ω sin λ (10.5.3)
ẋ ẏ 0
= (−ẏω sin λ) êx + (ẋω sin λ) êy + (−ẏω cos λ) êz (10.5.4)
i i
i i
i i
ẍ + β 2 x =2ωz ẏ
)
(10.5.6)
ÿ + β 2 y = − 2ωz ẋ
where β 2 = g/`, and ωz = ω sin λ is the êz component of the angular velocity. Notice that
(10.5.6) consists of a pair of coupled second order differential equations. The equations
are called coupled because the equation for ẍ contains a term with ẏ and vice-versa. A
common √ means of solving coupled ODEs analytically is by multiplying one of the equations
by i = −1 and adding the two equations together. The result of multiplying the ÿ equation
by i and adding it to the ẍ equation is,
q̈ + 2iωz q̇ + β 2 q = 0 (10.5.8)
The form of (10.5.8) should be familiar; it is similar to a damped harmonic oscillator.
The solution to (10.5.8) is,
q q
−iωz t
q(t) = e A exp it ωz + β + B exp −it ωz + β
2 2 2 2 (10.5.9)
Notice that if the Earth were not rotating, ωz = 0 and (10.5.8) would become,
q̈ 0 + β 2 q 0 = 0 (10.5.10)
which is the equation for simple harmonic motion, where we used the notation of q 0 for the
value of q in a non-rotating frame. From our work in Chapter 6, we know that
x + iy =(x0 + iy 0 )e−iωz t
= x0 + iy 0 (cos (ωz t) − i sin (ωz t)) (10.5.14)
Next, we equate the real and imaginary parts to get the solution we are after, x(t) and
y(t),
)
x(t) =x0 cos (ωz t) + y 0 sin (ωz t)
(10.5.15)
y(t) = − x0 sin (ωz t) + y 0 cos (ωz t)
i i
i i
i i
ez (vertical)
v0
ey
(East)
ex
(South)
Figure 10.8: The coordinate system for a projectile launched in the yz-plane towards the
East, at an angle α with a speed v0 .
It is easier to get some physical insight into these equations if we rewrite (10.5.15) as a
matrix equation,
a = g − 2 (ω × v) (10.6.1)
i i
i i
i i
Other than gravity, there are no external forces acting on the system. Therefore, F0 = 0.
As with the Foucault pendulum, the effective acceleration due to gravity is g = −gêz . The
Coriolis acceleration is,
êx êy êz
ω × v = −ω cos λ 0 ω sin λ (10.6.2)
ẋ ẏ ż
ω × v = (−ẏω sin λ) êx + (żω cos λ + ẋω sin λ) êy + (−ẏω cos λ) êz (10.6.3)
Notice that we made no assumptions about the components of the projectile’s velocity.
Inserting the accelerations into (10.6.1), we obtain the following equations,
ẍ =2ẏω sin λ
ÿ = − 2ω (ż cos λ + ẋ sin λ) (10.6.4)
z̈ = − g + 2ẏω cos λ
The result is a system of coupled ODEs. While it is possible to argue some approxi-
mations to simplify (10.6.4) (such as 2ẏω cos λ g), we will seek a numerical solution of
(10.6.4) using Mathematica. Note that statements between (* and *) are comments.
i i
i i
i i
long time and have Mathematica find the root of z(t) using the FindRoot command. Once
we have the time of flight T , we need to evaluate the numerical solution at t = T . We find
that the deflection is about 6.5 m. The range of this projectile is around 4080 m, so the
deflection is small but significant if one needs high precision.
ma = F0 − mA − 2m (ω
ω × v) − m (ω̇
ω × r) − m [ω
ω × (ω
ω × r)]
where:
a = the acceleration relative to the noninertial frame
F0 = the sum of the “interaction” forces acting on the particle (gravity,
friction)
A = the linear acceleration of the noninertial frame relative to an inertial one
−2 (ωω × v) = the Coriolis acceleration
ω × r = an acceleration due to the angular acceleration of the noninertial frame
−ω̇
ω × (ω
−ω ω × r) = the centrifugal acceleration.
The last four terms in the equation above are referred to as “inertial forces” or fictitious
forces, because they arise from the acceleration of the noninertial frame and are not due
to the interaction between the particle and another body.
i i
i i
i i
F0 = mA + ma + 2m (ω
ω × v) + m (ω̇
ω × r) + mω
ω × (ω
ω × r)
i i
i i
i i
q̈ + 2iωz q̇ + β 2 q = 0.
15. What is the rate of rotation of Foucault pendulum’s plane of oscillation located at the
equator? How about the North Pole? Explain your answers.
16. What is the rate of rotation of a Foucault pendulum’s plane of oscillation if the
pendulum is located in Paris, France? How about Paris, Kentucky, USA. Plot the
Foucault pendulum’s rotation rate as a function of latitude. Where is its maximum?
What is the minimum value? Where does the minimum value occur?
17. Consider a Foucault pendulum experiencing a damping force F0 = −bvv̂. Is the pre-
cession frequency changed?
18. Is the Foucault pendulum’s precession frequency significantly changed if ω̇ ω 6= 0, where
ω is the angular velocity of the Earth about its rotation axis? To answer this question,
assume that the magnitude of the Earth’s angular velocity changes at a constant rate,
but its direction does not change (this is not true over long time periods). Furthermore,
assume that the change in the angular velocity is very small. For example, consider
the situation described in Problem 26 from Chapter 4 which states that the Earth’s
angular velocity changed from 7.6 × 10−5 rad/s to 7.3 × 10−5 rad/s during a time
period of 350 million years. Let the pendulum have a mass m and a length ` (assumed
to be very long) and to be located at a latitude λ in the Northern Hemisphere.
Section 10.6: Projectile Motion in a Noninertial Frame
19. Make the approximation suggested in the text after (10.6.4), then solve the resulting
system of two coupled ODEs in closed form. You will want to use a symbolic ODE
solver such as Mathematica’s DSolve. After finding the solutions, perform a series
expansion for x(t) and y(t) and interpret your results.
20. Using the same initial velocity and latitude used in Section 10.6, find the range and
Coriolis deflection of a 1.0 kg projectile experiencing linear air resistance of the form
F0 = −bvv̂ where b = 0.3 kg/s.
21. Repeat Problem 20, but this time, the projectile experiences quadratic air resistance
of the form F0 = −bv 2 v̂ where b = 0.3 kg/m.
22. A projectile is fired due east from a point on the Earth’s surface with a latitude of
λ (in the Northern Hemisphere). The projectile’s initial speed is v0 and is launched
at an angle α with respect to the horizontal. Calculate the lateral deflection of the
projectile as it hits the ground. Do this problem without the aid of a computer.
23. There is a legend that during World War I, the British navy consistently missed
German ships when fighting near the Falkland Islands because their ships did not
properly account for the Coriolis force. We won’t discuss the validity of this legend
i i
i i
i i
here, but for a moment, let’s suppose that it is true and that the Coriolis force was
known, but the ship’s guns were set up to hit their targets for battles in the Northern
Hemisphere. The Falkland Islands are near 50◦ South latitude. Suppose the guns were
set to accurately hit targets at 50◦ North latitude, by how much did the British ships
miss their target during the Falkland Islands engagement? Assume that the German
ships were due east of the British ships and that the British ships’ guns had a muzzle
velocity of 500 m/s and the shells were shot at an angle of 20◦ with respect to the
horizontal.
i i
i i
i i
CHAPTER 11
We begin this chapter by reviewing the rotational motion of a single particle around an
arbitrary axis, and the concepts of the moment of inertia and the center of mass. In par-
ticular we focus on how the center of mass simplifies the description of the translational
and rotational motion of a system of particles. After the review, we explore generalized
definitions of the moment of inertia, including products of inertia and the inertia tensor,
and we demonstrate how to calculate these quantities for a variety of solids. We will see
that the moment of inertia tensor of a solid depends on the choice of the coordinate system.
Furthermore, we will discuss the parallel axis theorem for rigid bodies and show how it
can be used to calculate the moment of inertia tensor. This is followed by a discussion of
eigenvalues and eigenvectors of matrices, and how they can be used to describe the principal
axes of a rigid body. The chapter will conclude with a discussion of the Euler equations and
how they can be used to describe the precessional motion of spinning tops and gyroscopes.
335
i i
i i
i i
Figure 11.1: Rotational motion of a particle in the xy-plane, showing the position, velocity,
angular momentum, and angular velocity vectors.
Recall that unbolded variables represent magnitudes of vectors, e.g., ω = |ω ω |. Notice that
in this equation, the angular momentum is written as a quantity mr2 multiplied by the
velocity v. If we relate this to the definition of linear momentum p = mv, we can think of
(11.1.3) as a “rotational inertia” (mr2 ) multiplied by the angular velocity ω. This rotational
inertia is called the moment of inertia I, and is defined by:
I = mr2 (11.1.4)
where r is the distance of the mass m from the rotational axis.
1
The kinetic energy of the rotating mass m is found from T = mv 2 , and by substituting
2
v = ωr we obtain:
1 2
T = Iω (11.1.5)
2
Recall from Chapter 4 that the torque acting on the particle is N = r × F. However, we
can rewrite the torque as the time derivative of the angular momentum, by using Newton’s
second law for rotational motion. Using (11.1.3) we can write ` = Iω , and Newton’s second
law of rotation becomes:
d
N= (Iω) = I ω̇ = Iα (11.1.6)
dt
where α = ω̇ is the angular acceleration.
Let’s reconsider (11.1.3) as a vector relationship:
` = Iω
ω (11.1.7)
Since the moment of inertia is defined as a scalar quantity in (11.1.4), this vector equation
states that in the example of Figure 11.1, the angular momentum vector ` and the angular
velocity vector ω point in the same direction along the rotational axis. In addition, we can
rewrite (11.1.6) as
i i
i i
i i
N = Iα
α (11.1.8)
note that the torque vector N and the angular acceleration vector α point in the same
direction in the case illustrated in Figure 11.1.
Example 11.1 shows an evaluation of the angular momentum vector using Mathematica.
Example 11.1: Example of angular momentum and angular velocity for Figure
11.1
Evaluate the angular momentum and torque N acting on the mass m in the situation
shown in Figure 11.1. Assume that the mass m moves with a constant angular velocity Ω
around the z-axis.
Solution:
The Mathematica code is shown in Algorithm 10. We use cylindrical coordinates with
the angle increasing with time according to θ = Ωt. The program calculates the cross
product ω × r using the Cross command and verifies that this is indeed equal to the
velocity vector (dx/dt, dy/dt).
In this example, the angular momentum vector points always along the z-axis and its
magnitude is equal to `z = mR2 Ω.
The torque N = I Ω̇ = 0, since the mass moves with a constant angular velocity Ω
around the z-axis.
x = R ∗ Cos[Ω ∗ t];
y = R ∗ Sin[Ω ∗ t];
z = 0;
r = {x, y, z};
v = {D[x, t], D[y, t], D[z, t]}
ω = Ω ∗ {0, 0, 1};
Cross[ω, r] == v
OUTPUT: True
L = m ∗ Cross[r, v]//Simplify
OUTPUT: 0, 0, mR2 Ω
Let us now consider a particle of mass m which rotates around the z-axis at a constant
angle θ and with an instantaneous angular velocity ω as shown in Figure 11.2. In this
situation, the angular momentum vector is perpendicular to the shaded plane in the figure,
which is defined by the position and velocity vectors. As a result, in the example of Figure
11.2 the angular momentum vector ` and the angular velocity vector ω do not point in the
same direction, and the direction of ` changes continuously in space during the rotational
motion.
i i
i i
i i
Figure 11.2: A particle of mass m revolves around the z-axis on a plane parallel to the
xy-plane with angular velocity vector ω. In this example, the angular momentum vector `
and angular velocity vectors ω are not pointing in the same direction. Compare this figure
with the situation in Figure 11.1.
In this situation the equation ` = Iω does not apply, because ` and ω do not point in
the same direction in Figure 11.1. We must redefine the moment of inertia concept for the
mass m. As we will see later in this chapter, we need to replace the scalar quantity I in this
equation with a more general physical quantity, the moment of inertia tensor I.
Example 11.2: Example of angular momentum and torque for Figure 11.2
Evaluate the angular momentum and torque for the rotating mass m shown in Figure
11.2, assuming that ω = |ω
ω | is constant.
Solution:
We follow the same method as in the previous example, using cylindrical coordinates.
The Mathematica code appears in Algorithm 11. Clearly in this situation the angular
momentum vector ` has nonzero components along the x-, y-, and z-axes. By inspection
of the calculated components `x and `y , we see that the vector ` rotates around the z-axis
with a constant angular speed Ω, so that the z-component `z = mR2 Ω sin2 θ stays constant
with time.
Before we discuss the rotational properties of solid bodies, let us review the center of
mass, and how it is used to describe the rotation of a system of particles.
i i
i i
i i
ω = Ω ∗ {0, 0, 1};
Cross[ω, r] == v
OUTPUT: True
L = m ∗ Cross[r, v]//Simplify
ri'=ri-R
mi
R
ri
x
Figure 11.3: A collection of discrete particles (filled circles), with the empty circle repre-
senting their center of mass located at R.
Let R represent the location of the center of mass relative to the origin, shown as a
white circle in Figure 11.3, and r0i denote the positions of the particles relative to the center
of mass. Then ri = r0i + R and the position of the center of mass can be found from:
i i
i i
i i
Ṗ = M R̈ = Fext (11.2.4)
As a consequence, if the total external force acting on a system of particles is zero, then
the center of mass is either at rest, or it moves with constant velocity and the total linear
momentum of the system is conserved.
i=1
L̇ = N (11.2.6)
As a consequence, if the total external torque acting on a system of particles is zero, then
the total angular momentum of the system is conserved.
i i
i i
i i
W12 = T2 − T1 (11.2.8)
Furthermore, for a rigid body, the only motion relative to the center of mass is rotation.
Hence, the second term in (11.2.7) is the kinetic energy of rotation about the center of
mass. We can then think of equation (11.2.7) as stating that the total kinetic energy of a
rigid body is equal to the translational kinetic energy of the center of mass plus the kinetic
energy of rotation about the rigid body’s center of mass.
For the potential energy, we need to consider both external and internal forces, and their
corresponding potential energies. The potential energy of the ith mass is:
X
Vi = Viext + int
Vi,j (11.2.9)
j6=i
where Viext is the potential energy due to the external forces acting on mi , and Vi,j int is the
potential energy due to the internal interaction force between mi and mj . An example of
such an interaction force could be gravitational, electrostatic, or any other central force. In
int = V int r − r , i.e., the potential energy depends only on
general for central forces, Vi,j
i,j i j
the distance between the two particles. Therefore the total potential energy of the system
is:
XN N X
X
V = Viext + int
(11.2.10)
Vi,j ri − rj
i=1 i=1 j6=i
Rigid bodies are defined as solids in which the particles that make up the solid are at
fixed distances from each other. Therefore, the internal potential energy can be ignored
since it will be constant for each particle. Hence, for a rigid body we will only have to
worry about external forces; the internal central forces are irrelevant. When all internal and
external forces acting on a system of particles are conservative, the total mechanical energy
E = T + V is conserved.
In Chapters 4 and 5 we saw that the above results for linear momentum, angular
momentum, and kinetic energy hold for continuous mass distributions with the summa-
tions replaced by appropriate integrals.
i i
i i
i i
We describe this rigid body as a collection of finite masses mi located at positions ri , with
respect to a coordinate system xyz which is fixed on the rotating body. The instantaneous
velocity vi of mass mi is given by the cross product of the angular velocity vector ω =
ωx î + ωy ĵ + ωz k̂ and the positions ri = xi î + yi ĵ + zi k̂:
vi = ω × ri (11.3.1)
Figure 11.4: A rigid body with a fixed axis of rotation. The coordinate system shown is
fixed on the rotating body.
The total angular momentum vector L relative to the center of the coordinate system
is given by the vector sum of the angular momenta of the masses mi :
X X
L= mi (ri × vi ) = mi [ri × (ω × ri )] (11.3.2)
where the summation i = 1 . . . n runs over the masses mi which make up the rigid body.
We now use the identity of the triple cross product A × (B × A) = A2 B − (A · B) A to
evaluate ri × (ω × ri ):
ri × (ω × ri ) = ri2 ω − ri (ri · ω) (11.3.3)
ri × (ω × ri ) = x2i + yi2 + zi2 ωx î + ωy ĵ + ωz k̂ − xi î + yi ĵ + zi k̂ (xi ωx + yi ωy + zi ωz )
(11.3.4)
Substituting (11.3.4) into (11.3.2) and collecting the terms with î, ĵ, and k̂:
L = Lx î + Ly ĵ + Lz k̂ (11.3.5)
i i
i i
i i
N N N
! ! !
X X X
Ixy = Iyx = − mi xi yi Ixz = Izx = − mi xi zi Iyz = Izy = − mi yi zi
i=1 i=1 i=1
(11.3.8)
From these equations it is clear that the moment of inertia matrix I is symmetric, i.e.,
Ixy = Iyx , Ixz = Izx , and Iyz = Izy .
The terms Ixx , Iyy , and Izz are referred to as moments of inertia. These are the moments
of inertia of the rigid body about each axis. The terms Ixy , Iyz , Ixz , . . . are called products of
inertia. As we will see later, the products of inertia are zero if the x, y, and z axes correspond
to the rigid body’s axes of symmetry. The products of inertia measure the symmetry of the
rigid body’s mass distribution about the axes, x, y, and z. To understand the importance
of the product of inertia, consider the following. Suppose you are interested in balancing an
automobile tire. In the case of the tire, the axis of symmetry is the axle passing through the
center of the tire and perpendicular to the plane of the tire. A proper tire rotates about the
axle, or in other words, its angular momentum vector points along the direction of the axle.
Suppose the axle points in the z-direction. We want L = Izz ω k̂ to be the angular momentum
of the wheel. However, suppose the mass of the wheel is not evenly distributed about the
axle, then the axle is no longer an axis of symmetry and there will be at least two nonzero
products of inertia. For simplicity, suppose that Izy 6= 0 and all other products of inertia
are equal to zero. Then in this case, the angular momentum of the wheel is
where we use the compact notation uis (s = 1, 2, 3) to denote the three components of the
position vector ri = (xi , yi , zi ) of the mass mi . For example, u32 is the y-coordinate of
m3 . The symbol δkl in (11.3.10) is the Kronecker delta, which equals 1 when k = l and 0
otherwise.
Equation (11.3.6) can be written as the product of two matrices I and ω:
i i
i i
i i
L = I·ω (11.3.11)
Ixx Ixy Ixz ωx
I = Iyx Iyy Iyz and ω = ω y (11.3.12)
Izx Izy Izz ωz
It is important to remember that the elements of the inertia tensor depend on the choice of
origin of the coordinate system. Example 11.3 shows how to calculate the elements of the
moment of inertia tensor for a single particle.
Solution:
The instantaneous position of the mass m is best described by cylindrical coordinates,
with the constant radius r = R sin φ, i.e.,
Similarly
Ixy = −mxy = −mR2 sin2 φ sin θ cos θ Ixz = −mxz = −mR2 sin φ cos φ cos θ
Iyz = −mxz = −mR2 sin φ cos φ sin θ
i i
i i
i i
The angular momentum vector can now be found by multiplying this matrix I with
the angular momentum matrix ω = (0, 0, ω).
` = Iω
After the matrix multiplication and using the trig identities cos2 φ − sin2 φ = cos (2φ)
and 2 cos φ sin φ = sin (2φ), we obtain:
`x = −mR2 ω sin (2φ) cos (ωt) `y = −mR2 ω sin (2φ) sin (ωt) `z = −mR2 ω cos (2φ)
The physical interpretation of these equations is that as the mass m rotates around
the z-axis at the fixed angle φ, the angular momentum vector ` rotates around the z-axis
also, but with a different angle equal to 2φ. The z-component of the angular momentum
vector stays constant in time, and is equal to `z = −mR2 ω cos (2φ).
If we are dealing with a continuous uniform distribution of particles instead of a collection
of discrete masses mi , the summations in the above equations become integrals over the
continuous variables (x, y, z), and the density of the material ρ(x, y, z) must be included
inside the integrals. For solids in three dimensions, the components of the inertia tensor
become:
Inertia Tensor for a Continuous Mass Distribution
ZZZ
Ixx = y 2 + z 2 ρdxdydz
V
ZZZ
Iyy = x + z ρdxdydz
2 2
(11.3.13)
V
ZZZ
Izz = x + y ρdxdydz
2 2
V
Example 11.4 shows how to evaluate the elements of the inertia tensor for a triangular
pyramid with a variable density, by using numerical integration in Mathematica.
Example 11.4: The moment of inertia of a triangular pyramid
Consider the triangular pyramid with vertices at the origin, (1, 0, 0), (0,1, 0), and (0, 0, 1)
shown in the figure below. The pyramid has a density ρ = y 2 ex sin x2 (note that this
density is not physical since it can have the value of zero). Find the moment of inertia
tensor with respect to the Cartesian coordinate system with the origin at (0, 0, 0).
i i
i i
i i
(0,0,1)
dV
y
(0,1,0)
(1,0,0)
Solution:
In order to find the moment of inertia tensor, we first need to choose a coordinate
system and identify the limits of integration. We will use Cartesian coordinates for the
calculation because it will be easiest to describe the pyramid in Cartesian coordinates.
Next, we need to figure out the limits of integration that we will need to evaluate (11.3.13)
and (11.3.14).
The volume element dV appears in the figure. Notice that its height is determined by
its x and y coordinates. The “front” face of the pyramid in the figure is described by the
equation x + y + z = 1. Therefore, the limits for the z integral will range from 0 to 1 − x − y.
Next, the distance the volume element can be translated along the y-axis in the xy-plane
is determined by its x-coordinate. Therefore, the limits for the y integral will be from 0
to 1 − x (the equation of the line that forms the base of the pyramid in the xy-plane).
Finally, x can range from 0 to 1.
Now that we have all of the pieces, we can perform the integrals. We will evaluate the
integrals in Mathematica using the algorithm below. Notice that the limits of the integral
that we would do last by hand, are the first limits in the NIntegrate command.
Rho[x , y , z ]:=y ∧ 2*Exp[x]*Sin[x];
tensor = {{Ixx, Ixy, Ixz}, {Ixy, Iyy, Iyz}, {Ixz, Iyz, Izz}};
MatrixForm[tensor]
i i
i i
i i
1 1 1 1
T = ω · L= Ixx ωx2 + Iyy ωy2 + Izz ωz2 +Ixy ω x ω y +Iyz ω y ω z +Izx ω z ω x (11.4.3)
2 2 2 2
By using (11.3.11), this equation can be written in a compact form as a matrix equa-
tion:
Kinetic Energy of a Rigid Body
1 1
T = (ω)T · L = (ω)T · I · ω (11.4.4)
2 2
Here the dot product indicates the multiplication of the 3 matrices, and (ω)T indicates the
transpose of the column matrix representing the angular velocity vector ω.
If the angular velocity vector ω points in the direction of the unit vector n̂, we can write
ω = ωn̂, where ω is the magnitude of the angular frequency. Therefore the above matrix
equation becomes:
1
T = ω 2 (n̂)T · I · n̂ (11.4.5)
2
From introductory physics we recall that the kinetic energy for rotation around a fixed axis
is:
1
T = In ω 2 (11.4.6)
2
where In is the scalar moment of inertia around this axis. By comparing the last two
equations, we obtain the general expression:
In = (n̂)T · I · n̂ (11.4.7)
The products of inertia can be found using (11.4.7). For example, Ixy is obtained by calcu-
lating the quantity:
T
Ixy = bi · I · bj (11.4.8)
i i
i i
i i
a2 m
OUTPUT: 6
2a2 m
OUTPUT: 3
5a2 m
OUTPUT: 12
i i
i i
i i
Figure 11.5: Two coordinate systems with their axes parallel to each other, are shifted in
space relative to each other by a vector a.
In its simplest form, the parallel axis theorem states that if I is the scalar moment of
inertia of a body around an axis AB, and ICM is the corresponding moment of inertia about
a second axis parallel to AB and passing through the center of mass, then I and ICM are
related by:
Iz = Ix + Iy (11.5.2)
More generally, it is convenient and useful to know the relationship between inertia tensors
expressed in different coordinate systems.
We now prove the parallel axis theorem. Figure 11.5 shows two Cartesian coordinate
systems with origins O and OCM and with their axes parallel to each other. Note that OCM
is located at the object’s center of mass. If a is the vector connecting the two origins, the
relationship of any position vectors ri and r0i in the two coordinate system is:
ri = r0i + a (11.5.3)
Let I be the inertia tensor defined in the coordinate system (x, y, z) with the origin fixed at
point O, and I0 be the inertia tensor for the center of mass coordinate system (x0 , y 0 , z 0 )with
its origin at OCM . We wish to find the relation between I and I0 . As we saw previously
i i
i i
i i
in (11.3.10), the components of the inertia tensor I with respect to O can be written in
compact form as:
N 3
!
X X
Ikl = mi δkl 2
uis − uik uil (11.5.4)
i=1 s=1
where we again use the compact notation uis (s = 1, 2, 3) to denote the three components
of the mass mi ’s position vector ri = (xi , yi , zi ), and δkl is the Kronecker delta.
Similarly, the components for the tensor with respect to the second coordinate system
OCM are:
N 3
!
0 2
X X
0 0 0
Ikl = (11.5.5)
mi δkl uis − uik uil
i=1 s=1
Substituting (11.5.6) into (11.5.4), by expanding the terms and rearranging, we obtain:
N 3
!
X X 2
Ikl = u0is + as − u0ik + ak u0il + al
mi δkl
i=1 s=1
N 3 N 3
" # " #
X X 2 X X
= mi δkl u0is − u0ik u0il + mi δkl a2s − ak al
i=1 s=1 i=1 s=1
N
X N N
X X
+2 mi δkl u0is as − mi u0ik ai − mi u0il ak (11.5.7)
The total mass is M = mi , so (11.5.10) gives the generalized parallel axes theorem, also
P
known as the Steiner theorem, for the elements of the inertia tensor in any coordinate
system:
i i
i i
i i
Note that I0 is the object’s moment of inertia for a coordinate system whose origin is
located at the object’s center of mass.
By applying the parallel axis theorem for the diagonal elements of the tensor so that
δkl = 1, we recover (11.5.1), the simple form of the parallel axis theorem:
0 0
Ikk = Ikk + M a2 − a2k = Ikk + M d2k (11.5.12)
where dk is the shortest distance from the axis of rotation to the center of mass.
λ1 0 0
I = 0 λ2 0 (11.6.1)
0 0 λ3
This special set of axes is known as the principal axes of the rigid body. According to the
same theorem from linear algebra, there are three principal axes of rotation and they are
mutually orthogonal. In addition, any axis of symmetry through the origin is a principal
axis.
The main property of the principal axes is that when rotation takes place around one
of these three axes, the angular momentum L and angular velocity vectors ω are parallel
to each other. Therefore, if ω points along the direction of a principal axis, then L = λωω
(where λ is a scalar), in order for L and ω to point in the same direction. However, recall
that in general, L = I · ω . Hence:
I · ω = λω (11.6.2)
The scalar λ is called a principal moment of inertia, and it is the moment of inertia of the
object when it is rotating about the principal axis ω . Before discussing principal moments
and principal axes further, it is important to pause and discuss some general properties of
(11.6.2).
The equation I · ω = λω ω is an eigenvalue problem in the theory of linear algebra. In
general, in an eigenvalue problem we are given a square matrix B, and we are looking for
two quantities, a vector a and a scalar quantity λ, such that:
Ba = λa (11.6.3)
The vector a is called an eigenvector of the square matrix B, corresponding to the eigenvalue
λ. Hence, according to (11.6.2), the principal axes are the eigenvectors of the inertia tensor,
and principal moments are the eigenvalues associated with each principal axis.We see from
(11.6.3), that the matrix B scales the vector a by a value λ.
i i
i i
i i
Physically, from (11.6.3) we see that for λ > 0, the matrix B elongates or shortens the
vector a, but does not change the direction of a. If λ = 1, then a remains unchanged after
multiplication by B. When λ < 0, the vector a is still scaled, but now it points in the
direction opposite of its original orientation.
According to another theorem in linear algebra, in order for the eigenvalue equation
(11.6.3) to have a nontrivial solution a 6= 0, the determinant of the matrix (B − λ1) must
be zero, i.e.:
det (B − λ1) = 0 (11.6.4)
where 1 is the square identity matrix, consisting of ones along the diagonal and zeros
everywhere else. Equation (11.6.4) is called the characteristic equation of our eigenvalue
problem. In Chapter 13, we will revisit the eigenvalue problem in more detail.
There are always two steps in finding the eigenvalues and eigenvectors of a given moment
of inertia matrix:
1. Solve the characteristic equation (11.6.4) in the form det (I − λ1) = 0, in order to
find the 3 eigenvalues λ1 , λ2 , and λ3 . These eigenvalues will depend of course on the
physical parameters of the rigid body.
2. Substitute the first eigenvalue λ1 into Ia1 = λ1 a1 , in order to find the corresponding
eigenvector, a1 . Next, repeat this step for each of the remaining eigenvalues, in order
to obtain their eigenvector.
At the end of this process, we will have found three eigenvectors a1 , a2 , and a3 and the
corresponding eigenvalues λ1 , λ2 , and λ3 , respectively. The resulting eigenvalues are the
principal moments of inertia and the resulting eigenvectors are the principal axes. The
principal moment of inertia λi is the rigid body’s moment of inertia for a rotation about
the principal axis ai .
The Mathematica code in Examples 11.6 and 11.7 shows how to evaluate the moment
of inertia tensor and moments of inertia of a cube, with respect to two different coordinate
systems.
i i
i i
i i
ρ = m/a∧ 3;
Ixx = Iyy = Izz = Integrate[ρ ∗ (y ∧ 2 + z ∧ 2), {x, 0, a}, {y, 0, a}, {z, 0, a}]
2a2 m
OUTPUT: 3
Ixy = Iyz = Ixz = −Integrate[ρ ∗ (x ∗ y), {x, 0, a}, {y, 0, a}, {z, 0, a}]
2
OUTPUT: − a 4m
Eigenvectors[iCube]
vec1 = Eigenvectors[iCube][[1]];
vec2 = Eigenvectors[iCube][[2]];
vec3 = Eigenvectors[iCube][[3]];
i i
i i
i i
a2 m
OUTPUT: 6
Ixy = Iyz = Ixz = −Integrate[ρ ∗ (x ∗ y), {x, −a/2, a/2}, {y, −a/2, a/2}, {z, −a/2, a/2}]
OUTPUT: 0
Eigenvectors[iCube]
vec1 = Eigenvectors[iCube][[1]];
vec2 = Eigenvectors[iCube][[2]];
vec3 = Eigenvectors[iCube][[3]];
i i
i i
i i
dL⟂
r CM
mg
x
y
Figure 11.6: The spinning top precesses at a fixed angle θ due to the gravitational torque.
The top is rotating about its symmetry axis which points in the same direction as L.
i i
i i
i i
Let us take the contact point to be the origin O of the coordinate axes. The total
angular momentum of the spinning top with respect to this origin is the sum of the angular
momentum Lspin (due to the spinning motion) and the angular momentum LCM (due to
the motion of the center of mass about the pivot point):
i i
i i
i i
i i
i i
i i
Figure 11.7: The motion of a rotating rigid body rotating about the axis AB can be described
in two coordinate systems. The space coordinate frame X 0 Y 0 Z 0 is an inertial coordinate
system. The body coordinate frame XY Z is fixed on, and rotates with the rigid body, and
is therefore a non-inertial frame. The body coordinate frame usually corresponds to the
principal axes of the rotating rigid body.
Newton’s second law in an inertial frame of reference says that the time derivative of
the angular momentum L in the inertial frame equals the applied torque N:
dL d
= (Ispace ω) = N (11.8.1)
dt space dt
where Ispace is the moment of inertia tensor calculated in the space (inertial) frame, and ω =
(ω1 , ω2 , ω3 ) is the angular velocity. This equation does not lend itself easily to calculations
because both Ispace and ω change during the motion. However, the calculation can be
simplified using the body frame.
Let ê1 , ê2 , and ê3 be the unit vectors of the body frame pointing along the principal
axes of the rotating body, i.e., the axes X,Y, and Z in Figure 11.7. In the body frame, the
moment of inertia tensor is constant and diagonal:
I1 0 0
I = 0 I2 0 (11.8.2)
0 0 I3
where I1 , I2 ,and I3 are the object’s principal moments of inertia. The angular momentum
in the body frame L can be written as:
i i
i i
i i
By substituting L = Iω in this equation, and since the inertia tensor and the principal
moments Ik do not depend on time, we have:
dL
= Iω̇
ω (11.8.5)
dt body
Therefore, we arrive at the following vector form of Euler’s equations:
Euler’s Equations
where Nk (k = 1, 2, 3) are the components of the applied torque. This system of three cou-
pled differential equations for (ω1 , ω2 , ω3 ) is known as Euler’s equations for a rigid body. In
special cases they can be solved analytically, however in most cases they must be integrated
numerically, in order to obtain the time dependent functions (ω1 (t) , ω2 (t) , ω3 (t)).
i i
i i
i i
η̇ − iΩη = 0 (11.9.5)
where η = ω1 + iω2 . The solution to (11.9.5) is
η = AeiΩt (11.9.6)
Note that once we introduce complex numbers, the constant in the solution can also be
complex. In that case, we write A = αeiβ and (11.9.6) becomes
η = αei(Ωt+β) (11.9.7)
where α and β are real constants. The term β is a phase term, and it can be set equal to
zero by choosing an appropriate initial time for the problem. This is similar to observing a
sinusoidal wave as it passes by. By setting the initial time t0 = 0 to be the time when a crest
passes your position, the resulting wave motion can then be described as a cosine function
with the phase equal to zero. In what follows, we continue to use (11.9.6) as the solution to
(11.9.5), with A as a real constant.
Using the Euler relationship for exponentials, we can rewrite (11.9.6) as
i i
i i
i i
x'3
x3
L
'
ω ω x3
ϕ
x1 x'1
x2 x'2
Figure 11.8: The precession of ω about the symmetry axis of the force-free symmetric top
as observed in the (a) body frame and in the (b) space frame.
In addition to the precession due to the gravitational torques provided by the Sun and
the Moon, the Earth experiences another precession due to the Earth’s oblateness. The
Earth is slightly flattened at the poles due to its rotation (see Chapter 10). The result of
the flattening is that the Earth’s principal moments of inertia are not all identical. The
Earth’s principal moment of inertia about its polar axis is slightly larger than its other two
moments (by about 1 part in 300). Hence (11.9.9) tells us that there should be a precession
with a frequency of ω3 /300 where ω3 is equal to one rotation per day. The result is a small
precession (or wobble) of the Earth’s polar axis with a period of 300 days. The wobble was
discovered by an amateur astronomer Seth Chandler (1846 - 1913). However the period is
closer to 400 days because the Earth is not a perfectly rigid rotator.
i i
i i
i i
L = I·ω
Ixx Ixy Ixz ωx
I = Iyx Iyy Iyz and ω = ω y
Izx Izy Izz ωz
where the scalars Ikl (k, l = 1, 2, 3) are the components of the inertia tensor I, defined by:
X X X
Ixx = mi yi2 + zi2 = mi ri2 − x2i Ixy = Iyx = − etc.
mi xi yi
The elements of the inertia tensor depend on the choice of origin for the system. For
a mass m with density ρ these elements are given by:
ZZZ ZZZ
Ixx = ρ y + z dxdydz
2 2
Ixy = Iyx = − etc.
ρxydxdydz
V V
The kinetic energy and moment of inertia tensor are related by the matrix multiplication
equation:
1 1
T = ω· L = (ω)T ·I · ω
2 2
where (ω)T indicates the transpose of the column matrix ω.
The scalar moment of inertia I for rotation around an axis in the direction n̂ of any
unit vector, is found from the inertia tensor using the matrix multiplication equation:
In = (n̂)T · I · n̂
The parallel axis theorem for scalar moments of inertia states: if I is the scalar moment
of inertia of a body around an axis AB, and ICM is the corresponding moment of inertia
about a second axis parallel to AB and passing through the center of mass, then I and
ICM are related by:
I = ICM + md2
The generalized parallel axis theorem between moment of inertia tensors is as follows: If I
is the inertia tensor defined in a coordinate system with the origin fixed at point O, and
I0 is the inertia tensor defined in a center-of mass coordinate system with its origin at the
i i
i i
i i
center-of-mass OCM and whose axes are parallel to the previous coordinate system, then
the elements of the tensors I and I0 are related by:
0
Ikl = Ikl + M a2 δkl − ak al
L = I · ω = λω
2. Find the principal axes and principal moments of inertia for a system of two masses
m1 , m2 connected by a massless rigid rod of length l.
3. A cylinder has mass m and radius R, and is attached to a hanging mass m. The
cylinder rolls without slipping upwards on an inclined plane, as shown in Figure 11.9.
The string is wrapped around the cylinder.
(a) Find the magnitude and direction of the acceleration of the hanging mass.
(b) What are the magnitude and direction of the force of static friction at the contact
point between the cylinder and the inclined plane?
i i
i i
i i
Figure 11.9: Problem 11.3: A cylinder on an inclined plane attached to a hanging mass.
4. A uniform hoop of mass m and radius R hangs in a vertical plane and is supported
by a nail at a point on the circumference, as shown in Figure 11.10. Calculate the
natural frequency of small oscillations.
Figure 11.10: Problem 11.4: A uniform hoop supported by a nail, located at the circle.
5. Baseball players always try to hit what is known as the sweet spot on the baseball
bat. The sweet spot is where the hit delivers the most energy to the ball, while at
the same time minimizes the force on the hands of the player. To describe the physics
of the sweet spot, the baseball bat can be assumed to be rotating in space around a
pivot point at its end without linear displacement, as shown in Figure 11.11. At what
distance D along the baseball bat is the sweet spot located?
Figure 11.11: Problem 11.5: A baseball bat hitting a baseball at the sweet spot D away
from the pivot point.
i i
i i
i i
where I1 is the moment of inertia of the attached object around the rotational axis,
and A is the torsion coefficient describing the stiffness of the wire. Consider the torsion
pendulum shown in Figure 11.12, which consists of a thin disk of radius R and mass
M with a cylindrical mass M of radius R/4 placed on top of the disk. Find the ratio of
the periods of this torsion pendulum, with and without the presence of the cylindrical
mass on top of the thin disk.
7. A bowling ball of uniform density is thrown along a horizontal alley with initial velocity
vo in such a way that it initially slides without rolling. The ball has mass m, the
coefficient of static friction with the floor is µs , and the coefficient of sliding friction
with the floor is µk . Compute how far the ball will slide before it starts rolling. You
can ignore the effect of air resistance.
8. A coin is spinning about its axis of symmetry through its center with angular frequency
ω as shown in Figure 11.13A below. The coin is placed down on a horizontal surface
as shown in Figure 11.13B. The coin stops slipping and starts rolling away. What is
the velocity of the coin when it rolls away?
(A) (B)
9. A wheel of mass M and radius R is projected along a horizontal surface with an initial
linear velocity vo and an initial angular velocity ωo , as shown in Figure 11.14. The
wheel starts sliding along the surface. Let the coefficient of friction between the wheel
and the surface be µ.
i i
i i
i i
(a) How long does it take for the wheel to stop sliding?
(b) What is the velocity of the center of mass of the wheel when the slipping stops?
10. A solid uniform cylinder of mass m and radius R is placed on a plane inclined at angle
θ with respect to the horizontal, as shown in Figure 11.15. Let a be the acceleration
of the axis of the cylinder along the incline. The coefficient of friction between the
cylinder and plane is µ. For θ less than some critical angle θc , the cylinder will roll
down the incline without slipping.
(a) What is the critical angle θc ?
(b) What is the acceleration a for angles less than the critical angle?
11. A wheel of radius R and moment of inertia I is mounted on a frictionless axle through
its center. A flexible, weightless cord is wrapped around the rim of the wheel and
carries a body of mass M which begins descending as shown in Figure 11.16. What is
the tension in the cord?
i i
i i
i i
12. A thin uniform stick of mass m is resting with its bottom end on a frictionless table.
The stick is released from rest at an angle θo with respect to the vertical, as shown in
Figure 11.17. Find the force exerted by the table upon the stick at an infinitesimally
small time after its release.
Figure 11.17: Problem 11.12: A thin uniform stick in the process of falling over.
i i
i i
i i
Figure 11.18: Problem 11.14: A billiard ball being struck by a horiztonal cue.
Section 11.3: Angular Momentum, Angular Velocity, and the Inertia Tensor
15. A square plate with mass m and side a, has a uniform mass density and lies on the
xy-plane of a xyz coordinate system, with the origin located at one of its corners. The
square plate has a thickness b.
(a) Find the moment of inertia tensor of the square plate with respect to this coor-
dinate system.
(b) Find the principal moments of inertia and principal axes for this plate, when the
thickness b << a , i.e., when the plate is very thin.
16. Two equal point masses M are connected by a massless rigid rod of length 2b, to form
a dumbbell. The system is constrained to rotate about an axis fixed to the center of
the rod as shown in Figure 11.19. The angular velocity ω is a constant in time.
(a) Choose an appropriate coordinate system and calculate all elements of the inertia
tensor.
(b) Find and discuss the principal axes and principal moments of inertia.
Figure 11.19: Problem 11.16: Two point masses connected by a massless rod.
17. Find the principal moments of inertia and principal axes for a uniform cylinder of
height h and radius R.
Section 11.4: Kinetic Energy and the Inertia Tensor
18. Four masses, all of value m, lie in the xy-plane at positions (x, y) =
(a, 0), (−a, 0), (0, +2a), and (0, −2a). These are joined by massless rods to form
a rigid body.
i i
i i
i i
(a) Find the moment of inertia tensor using the xyz-axes as a reference system.
(b) Consider a direction given by a unit vector n̂ that makes equal angles with the
positive xyz-axes. Find the moment of inertia for rotation about this axis.
19. A compound pendulum has mass m and principal moments of inertia I1 , I2 , and I3 .
The pendulum oscillates about a horizontal axis which makes angles α, β, and γ with
respect to the principal axes of inertia. Show that the period of small oscillations is
r
mgd
T = 2π
I
where d is the distance from the center of mass to the axis of rotation, and
20. A spherical ball of mass m and radius r rolls without slipping on a track as shown in
Figure 11.20. Find the minimum height h above the top position in the loop that will
permit the ball to maintain constant contact with the rail of the loop.
21. A uniform plank of length 2a is held temporarily so that one end leans against a
frictionless vertical wall, and the other end rests on a frictionless floor making an
angle θ0 = 60° with the floor, as shown in Figure 11.21. When the plank is released,
it will slide down under the influence of gravity. Show that the time t that it will take
the plank to reach a new angle θ is:
Z t
dθ
t= q
3g
2a (sin θ0 − sin θ)
0
i i
i i
i i
Figure 11.21: Problem 11.21: A uniform plank temporarily leaning against a frictionless
vertical wall.
22. A uniform solid ball of radius a rolls with velocity v on a level surface, and collides
inelastically with a step of height h < a, as shown in Figure 11.22. Find the minimum
velocity for which the ball will “trip” up over the step. Assume that no slipping occurs
at the impact point.
23. A particle of mass m and speed v collides elastically with the end of a uniform thin
rod of mass M and length L as shown in Figure 11.23. The collision takes place on a
horizontal surface, so gravity can be ignored. Show that if mass m is stationary after
the collision, then M = 4m.
Figure 11.23: Problem 11.23: A particle colliding with the end of a uniform thin rod.
i i
i i
i i
24. Solve the following problems involving different physical pendulums. A physical pen-
dulum is a rigid body that swings freely about some pivot point.
(a) A rectangular thin plate with sides (a, b) is hung vertically
p from its edge of length
a. Show that the period of small oscillations is T = 2π 2b/(3g)
(b) A uniform solid sphere with radius R is hung vertically from p a point on its
surface. Show that the period of small oscillations is T = 2π 7a/(5g)
(c) In a simple pendulum, the string is replaced with a solid thin rod of mass M
and length L, and the mass m is attached at the end of this rod. Show that the
period of small oscillations is T = 2π 2(M + 3m)L/ [3(M + 2m)g])
p
25. Given that the moment of inertia of a cube about an axis that passes through the
center of mass and the center of one face is I0 , find the moment of inertia about an
axis passing through the center of mass and one corner of the cube.
1 0 0
M R2
I= 0 1 0
4
0 0 2
(b) Find the moment of inertia tensor of the combination of disk and point mass
when the system rotates around the pivot point A at the origin, in the coordinate
system shown in the figure.
(c) Find the principal moments and the principal axes.
Figure 11.24: Problem 11.26: A thin disk with a point mass (solid dot) attached at the edge.
i i
i i
i i
27. A torsion pendulum consists of a vertical wire attached to a cube, which may rotate
about the vertical. The cube is hung from a corner, then is hung from midway along
an edge, and finally it is hung from the middle of a face. Which one of the three
periods of oscillation of the cube is largest?
28. Consider a rotating object which is not experiencing any external torques. Using
Euler’s equations, show that the magnitude of the angular momentum is constant.
Also show that the rotational kinetic energy is constant.
29. A thin rectangular plate has dimensions a, and b. Find the torque required to rotate
this plate with constant angular velocity ω around one of its diagonals.
30. A uniform rigid wheel is located on the xy-plane, and has principal moments of inertia
I1 = I2 6= I3 about its body-fixed principal axes x̂1 , x̂2 and x̂3 as shown in Figure 11.25.
The wheel is attached at its center of mass to a bearing which allows frictionless
rotation about one space-fixed axis. The wheel is “dynamically balanced,” i.e., it can
rotate at constant ω, so that it can exert no torque on its bearing.
(a) Examine and integrate the Euler equations.
(b) What conditions must the components of ω satisfy for this dynamically balanced
system? Discuss the motion
Figure 11.25: Problem 11.30: A uniform rigid wheel rotating with along the axis ω .
i i
i i
i i
CHAPTER 12
Coupled Oscillations
In this chapter, we will explore the properties of coupled harmonic oscillators. These systems
can be analyzed by using either the Lagrangian approach of Chapter 8, or alternatively
using Newton’s second law. The simplest form of these systems in mechanics contains two
masses connected by springs to each other. A second simple example of coupled mechanical
oscillators is the double pendulum, which also exhibits a wide range of interesting behaviors.
We will see that these simple oscillating systems can exhibit normal modes of oscillation,
which are patterns of motion in which all parts of the system move sinusoidally with the
same frequency. The frequencies of the normal modes of a system are known as its natural
frequencies of oscillation. We will find that any motion exhibited by the system can be
expressed as a linear combination of these normal modes.
The discussion of the two-mass system will lead us to a more general description of
linearly coupled harmonic systems, and how their equations of motion can be written in
matrix form. The best way to obtain solutions to the equations of motions for coupled oscil-
lations is by using standard techniques from linear algebra, in order to find the eigenvalues
and eigenvectors of a matrix. The eigenvectors and eigenvalues of the matrix characterizing
the oscillating system are closely related to its normal modes.
This chapter will conclude with a general treatment of coupled oscillations and a dis-
cussion of normal coordinates.
Figure 12.1: A system of two coupled harmonic oscillators consisting of two masses m1 and
m2 connected with three springs with constants k1 , k2 , and k3 .
373
i i
i i
i i
Let us denote by x1 (t) and x2 (t) the horizontal displacements of the two masses from
their respective equilibrium points. The force on the first mass due to the first spring is
−k1 x1 . The middle spring will be stretched by a distance (x1 − x2 ), and the force on the
first mass due to this middle spring will be −k2 (x1 − x2 ). The total force on the first mass
must then be F1 = −k1 x1 − k2 (x1 − x2 ). Similarly, the force on the second mass due to the
middle spring is −k2 (x2 − x1 ), and the force on the second mass due to this third spring
will be −k3 x2 .
Using the notation ẍ for acceleration, the equations of motion from Newton’s second
law F = ma for the two masses are:
)
m1 ẍ1 = −k1 x1 − k2 (x1 − x2 )
(12.1.1)
m2 ẍ2 = −k2 (x2 − x1 ) − k3 x2
In a more general way, we can obtain the same equations by starting with the Lagrangian
formulation of Chapter 8. The potential energies of the two end springs are V1 = k1 x21 /2
and V3 = k3 x22 /2, while the potential energy for the middle spring is V2 = k2 (x1 − x2 )2 /2,
so that the Lagrangian is equal to:
1 1 1 1 1
L = T − VT otal = m1 x˙1 2 + m2 x˙2 2 − k1 x21 − k3 x22 − k2 (x1 − x2 )2 (12.1.2)
2 2 2 2 2
The Euler-Lagrange equations are:
d ∂L ∂L
− =0 → m1 ẍ1 = −k1 x1 − k2 (x1 − x2 ) (12.1.3)
dt ∂ ẋ1 ∂x1
d ∂L ∂L
− =0 → m2 ẍ2 = −k2 (x2 − x1 ) − k3 x2 (12.1.4)
dt ∂ ẋ2 ∂x2
These are of course the same equations as in (12.1.1).
In general, it is not possible to obtain the solutions x1 (t) and x2 (t) of the system of
equations (12.1.1) analytically, and they must be obtained by numerically integrating the
equations for given initial conditions of the system. The initial conditions are usually given
as the initial positions and initial speeds of the two masses.
Example 12.1 shows how to obtain and plot the numerical solutions x1 (t) and x2 (t),
with the initial conditions x1 (0) = 0, x2 (0) = a, ẋ1 (0) = a, and ẋ2 (0) = 0.
Example 12.1: Numerical solution for the general case of two coupled oscil-
lating masses
Integrate (12.1.1) for k1 = 1 N/m, k2 = 2 N/m, k3 = 3 N/m, m1 = 1 kg, m2 = 2 kg and
plot the numerical solutions x1 (t) and x2 (t), with the initial conditions x1 (0) = 0, x2 (0) =
1m, ẋ1 (0) = 0 and ẋ2 (0) = 0. This situation corresponds to the case where the first mass m1
is initially at rest at its equilibrium position (x1 (0) = 0 and ẋ1 (0) = 0), and the second mass
is pulled a distance a = 1 meter from its equilibrium and released from rest (ẋ2 (0) = 0).
Solution:
In this example we can use either the NDSolve command in Mathematica to numerically
solve the differential equations, or alternatively use the DSolve command to obtain the
analytical solutions x1 (t) and x2 (t). Although Mathematica can obtain analytical expres-
sions for x1 (t) and x2 (t), they are not listed here because they are algebraically very
complex, and therefore the output is suppressed in this example by using the semicolon
(;) at the end of the command line.
i i
i i
i i
The parameter numValues in the code contains the numerical values for the parameters
m1 , m2 , k1 , k2 , and k3 in the form of a rule (→). These numerical values are needed in
order to plot the solutions using the Plot and GraphicsGrid commands.
sol = DSolve[{m1 ∗ x1”[t] == −k1 ∗ x1[t] + k2 ∗ (x2[t] − x1[t]), m2 ∗ x2”[t] == −k3 ∗ x2[t] +
k2 ∗ (x1[t] − x2[t]), x1[0] == 0, x10 [0] == 0, x2[0] == a, x20 [0] == 0}, {x1[t], x2[t]}, t];
numValues = {a → 1, k1 → 1, k2 → 2, k3 → 3, m1 → 1, m2 → 2};
5 10 15 20 5 10 15 20
-0.5
-0.5
-1.0
-1.5 -1.0
The plots of x1 (t) and x2 (t) in Example 12.1 are obviously complex, and it is not
possible to give a simple physical description of the motion of the two masses. The key
physical component which creates this complex behavior is the middle spring in Figure
12.1, since this is the component that couples the motion of the two masses.
i i
i i
i i
and ẋ2 (0) = 0). The Mathematica code uses the DSolve command to obtain analytical
expression for x1 (t) and x2 (t), and the ExpToTrig command is used to convert the solutions
x1 (t) and x2 (t) into a form containing trigonometric functions, instead of exponential
functions. The result of DSolve appears above the graph. p
The analytical solutions are x1 (t) = x2 (t) = a cos k/m t . This tells us that if the
two masses are initially displaced from equilibrium by the same distance and released
from rest, the two masses will move together with the same speed and in phase as if the
middle spring was not present. This makes physical sense, since in this situation the middle
spring will be unstretched from its natural length, and will remain unstretched during the
motionp of the two masses. The frequency of oscillation for both masses in this situation is
ω1 = k/m, i.e., the same frequency as if only one of the two masses were attached to a
single spring with a spring constant k.
sol = ExpToTrig[DSolve[{m ∗ x1”[t] == −k ∗ x1[t] + k ∗ (x2[t] − x1[t]), m ∗ x2”[t] ==
−k ∗ x2[t] + k ∗ (x1[t] − x2[t]), x1[0] == a, x10 [0] == 0, x2[0] == a, x20 [0] == 0},
{x1[t], x2[t]}, t]]//Simplify
numValues = {a → 1, k → 1, m → 1};
OUTPUT:
nn h√ i h √ ioo
x1[t] → aCos √kt
m
, x2[t] → aCos √kt
m
x1[t] x2[t]
1.0 1.0
0.5 0.5
Out[ ]=
2 4 6 8 10 2 4 6 8 10
-0.5 -0.5
-1.0 -1.0
Example 12.3: Equal masses and identical springs: The antisymmetric oscil-
lation
Repeat Example 12.2, by using a different set of initial conditions x1 (0) = a, x2 (0) = −a,
ẋ1 (0) = 0 and ẋ2 (0) = 0. Plot the solutions with Mathematica by using the numerical values
a = 1.0 m, k = 1.0 N/m, m = 1.0 kg.
Solution:
In this situation, the two masses are initially displaced from their equilibrium positions
by equal and opposite distances a, and are then released from rest.
i i
i i
i i
The analytical solutions(appearing above the graph) in this case from Mathematica
are x1 (t) = a cos 3k/m t and x2 (t) = −a cos 3k/m t . This tells us that the two
p p
masses will move together with the same speed, but they will be completely out of phase
as shown in the output of the code. The frequency of oscillation for both masses in this
situation is ω2 = 3k/m.
p
numValues = {a → 1, k → 1, m → 1};
OUTPUT:
nn h√ √ i h √ √ ioo
x1[t] → aCos √3 kt
m
, x2[t] → −aCos 3 kt
√
m
x1[t] x2[t]
1.0 1.0
0.5 0.5
Out[ ]=
2 4 6 8 10 2 4 6 8 10
-0.5 -0.5
-1.0 -1.0
Examples 12.2 and 12.3 show that the system of two equal masses and three p
identical
springs in Figure 12.1 has two natural frequencies given by ω1 = k/m and ω2 = 3k/m.
p
By properly choosing the initial conditions in the system as in the two examples, we can
force both masses to oscillate with a single frequency, either ω1 or ω2 . In these special
situations, the two natural frequencies are uncoupled from each other, and we say that
these are the normal modes of the oscillating system.
Figure 12.2 shows schematically thep motion of the two masses in either the symmetric
oscillation pattern with frequency ω 1 = k/m (left panel), or an antisymmetric oscillation
with frequency ω2 = 3k/m (right panel).
p
i i
i i
i i
(A) (B)
Figure 12.2: The normal modes of the twopmass-three springpsystem in Figure 12.1, corre-
sponding to the natural frequencies ω1 = k/m and ω2 = 3k/m. These normal modes
represent symmetric and antisymmetric oscillations respectively.
and these are shown in panels (A) and (B) that follows. Once more the functions
x1 (t) and x2 (t) are complicated, and it is difficult to describe how exactly the two masses
are moving. This is because mathematically both x1 (t) and x2 (t) contain the frequencies
ω1 and ω2 . p
Panels (C) and (D) show plots of the function x1 (t) + x2 (t) = a cos k/m t and
p
x1 (t) − x2 (t) = a cos 3k/m t , respectively. By using these linear combinations, it is
now possible to decouple the two normal modes, so that the pmotions shown p in panels (C)
and (D) are simple cosine functions with frequencies ω1 = k/m and ω2 = 3k/m.
i i
i i
i i
numValues = {a → 1, k → 1, m → 1};
OUTPUT:
nn h√ i h √ √ i h√ i h √ √ ioo
x1[t] → 12 a Cos √kt
m
− Cos 3 kt
√
m
, x2[t] → 1
2 a Cos √kt + Cos
m
3 kt
√
m
1.0
0.5 0.5
2 4 6 8 10 2 4 6 8 10
-0.5 -0.5
-1.0 -1.0
1.0 1.0
0.5 0.5
2 4 6 8 10 2 4 6 8 10
-0.5 -0.5
-1.0 -1.0
In Example 12.5, we examine one more interesting behavior of the two-mass system,
the case of weakly coupled oscillators. In this example, the weak coupling is established by
choosing a middle spring with smaller spring constant than the two end springs (k1 = k3 = k
and k2 << k).
i i
i i
i i
i i
i i
i i
0.5 0.5
20 40 60 80 20 40 60 80
-0.5 -0.5
-1.0 -1.0
Out[ ]=
0.5 0.5
20 40 60 80 20 40 60 80
-0.5 -0.5
-1.0 -1.0
In the next two sections, we develop a more formal mathematical analysis of the normal
modes for the system in Figure 12.1 by using the techniques of linear algebra.
m 0
ẍ1 −2k k x1
= (12.2.2)
0 m ẍ2 k −2k x2
We can now solve this matrix equation by using the standard methods of linear algebra,
and by following the same eigenvalue problem procedure we used in Chapter 11 for the
principal moments of inertia.
We proceed in two steps, first we find the natural frequencies ω of the system, and second
we find the positions x1 (t) and x2 (t), as follows.
Since we expect oscillatory motion, we try solutions of the form:
x1 (t) = A1 eiωt x2 (t) = A2 eiωt (12.2.3)
i i
i i
i i
where A1 and A2 are the unknown amplitudes of oscillation for the two masses, and ω is
the unknown frequency of oscillation. Substituting these into the matrix equation (12.2.2):
m 0 −ω 2 A1 eiωt A1 eiωt
−2k k
= (12.2.4)
0 m −ω 2 A2 eiωt k −2k A2 eiωt
By canceling the exponential factor eiωt which is common to all terms, and combining the
matrices, we obtain:
−ω 2 m + 2k 0
−k A1
= (12.2.5)
−k −ω 2 m + 2k A2 0
This matrix equation represents a system of equations. A theorem from linear algebra says
that if the determinant of the matrix is nonzero, then there is a unique solution, which in
this case is the trivial solution A1 = A2 = 0. However, in order for multiple solutions to
exist, the determinant of the matrix must be zero. We are interested in a nontrivial solution
A1 , A2 6= 0, so we solve for the values of ω which cause the determinant to be zero.
We set the determinant of the matrix equal to zero:
−ω 2 m + 2k
−k
det =0 (12.2.6)
−k −ω 2 m + 2k
ω 2 m − 2k ω 2 m − 2k − k 2 = 0 (12.2.7)
Solving for ω, we obtain four possible solutions, only two of which are positive:
r r
k 3k
ω1 = , ω2 = (12.2.8)
m m
This completes the first part of the analysis, where we determined the two natural frequen-
cies ω1 and ω2 . In the previous section, we found that oscillatory solutions to our system of
equations can have one of these two frequencies.
In order to complete the description of the system, we mustq also find the two unknown
amplitudes of oscillation A1 and A2 . If we substitute ω1 = k
m into the matrix equation
(12.2.5), we obtain:
q 2
k
m m − 2k k
A1 −k k A1
A2 = k =0 (12.2.9)
q 2
−k A2
k k
m m − 2k
i i
i i
i i
x1 (t) 1
= D1 cos (ω1 t − φ1 ) (12.2.14)
x2 (t) 1
By working in a similar fashion for the second natural frequency of the system, we substitute
ω2 = 3k/m into (12.1.1), and obtain A1 = −A2 . Since A1 = −A2 , this type of motion
p
corresponds to the two masses moving in opposite directions, while the center of mass
remains stationary, as shown in Figure 12.2b. This type of motion is known as the second
normal mode or the antisymmetric mode of oscillation.
We can then write the second possible solution corresponding to ω2 = 3k/m as:
p
x1 (t) 1
= E1 cos (ω2 t − φ2 ) (12.2.15)
x2 (t) −1
A faster method of obtaining the normal mode frequencies ω1 and ω2 and the amplitudes
A1 and A2 , is by solving the matrix equation (12.2.5) using the symbolic capabilities of
Mathematica. Example 12.6 shows how to use the Solve command in Mathematica, to
obtain the general solution for the two-mass three-spring system.
Example 12.6: Solving the two-mass three-spring system by solving the matrix
equation
Solve the two oscillating mass system in Figure 12.1 as a matrix equation problem.
Find the normal mode frequencies ω1 and ω2 and the general relationship between the
amplitudes A1 and A2 , in the case of equal masses and identical springs.
Solution:
The Mathematica code below uses the Solve and Simplify commands to solve the matrix
equation (12.2.5). The Simplify command is used together with its option, Assumptions
which restricts the results to positive values of the parameters k and m. After a warning
message that Mathematica might not be able to obtain all solutions, the code produces
four possible frequencies, only two ofpwhich are positive and therefore acceptable. The
first normal mode frequency is ω1 = k/m, and the corresponding relationship between
the amplitudes is A1 = A2 . This is of course the symmetric mode of oscillation for the
systempthat we saw previously. The code also produces the second normal mode frequency
ω2 = 3k/m, and the corresponding relationship between the amplitudes A1 = −A2 for
the antisymmetric mode of oscillation.
i i
i i
i i
2 ∗ k/m
−k/m
M= ;
−k/m 2 ∗ k/m
A1 A1
Simplify Solve M. == ω ∧ 2 ∗ , {ω, A1, A2} ,
A2 A2
Assumptions->k > 0&&m > 0]
Solve : Equations may not give solutions for all “solve” variables.
OUTPUT:
q q
{A1 → 0, A2 → 0}, ω → − m k
, A2 → A1 , ω → m k
, A2 → A1 ,
√ qk √ qk
ω→− 3 m , A2 → −A1 , ω → 3 m , A2 → −A1
By substituting a trial solution of the form x1 (t) = A1 eiωt and x2 (t) = A2 eiωt and canceling
the common factor eiωt , these equations yield:
k1 k2 k2 k3
− A1 ω 2 = − A1 − (A1 − A2 ) − A2 ω 2 = − (A2 − A1 ) − A2 (12.2.18)
m1 m1 m2 m2
These can be written in compact matrix form as:
k1 +k2 −k2
m1 m1 A1 A1
= ω2 (12.2.19)
−k2 k2 +k3 A2 A2
m m 2 2
You will recognize that this equation is an eigenvalue problem in the theory of linear algebra,
similar to the eigenvalue problems we encountered for the moment of inertia matrix in
Chapter 11. As we remember from that chapter, in an eigenvalue problem we are given a
square matrix B, and we are asked to find a vector a such that Ba = λa, where λ is a
constant. The vector a is called an eigenvector of the square matrix B, corresponding to the
eigenvalue λ. We will explore the details of eigenvalue problems in Chapter 13.
For the two mass system of Figure 12.1, the eigenvalue problem to be solved becomes
clear by writing (12.2.19) in this matrix form:
i i
i i
i i
GA = ω 2 A (12.2.20)
k1 +k2 −k2
m1 m1
A1
A= G=
A2
−k2 k2 +k3
m2 m2
We are looking for the eigenvalues λ = ω 2 of the square matrix G, which will give us the
natural frequencies of oscillation. We are also looking for the corresponding eigenvectors A,
which will give us the normal modes of oscillation corresponding to each natural frequency.
In order for the eigenvalue equation(12.2.20) to have a nontrivial solution (A 6= 0), the
determinant of the matrix G − ω 2 1 must be zero, where 1 is the 2 × 2 identity matrix.
This gives:
k1 +k2 2 −k2
− ω
m1 m1
det =0 (12.2.21)
−k2 k2 +k3 2
m2 m2 − ω
This is the characteristic equation of our eigenvalue problem. From this point on, we proceed
by following the same two-step method used in the previous section. First, we must solve
the characteristic equation (12.2.21) in order to find the frequencies ω1 and ω2 . These
frequencies will depend on k1 , k2 , k3 , m1 , and m2 . In the next step, we substitute the first
natural frequency ω1 into (12.2.19), in order to find (A1 , A2 ), the first normal mode. Finally,
we repeat the previous steps using the second natural frequency ω2 , in order to find (A1 , A2 )
for the second normal mode.
The Mathematica and Python codes in Example 12.7 show how to find the eigenvectors
and eigenvalues for the two-mass system in Figure 12.1.
Example 12.7: The two-mass system as an eigenvalue/eigenvector problem
Solve the two oscillating mass system in Figure 12.1 as an eigenvalue/eigenvector prob-
lem, in these two cases:
(a) Identical springs and identical masses.
(b) Identical springs and different masses.
Solution:
(a) In this case the matrix,
2k −k
m m
G=
−k 2k
m m
and the Mathematica code below uses the commandspEigenvalues and Eigenvectors. The
code produces the first eigenvalue ω12 = k/m or ω1 = k/m, and the corresponding eigen-
vector,
1
A1
=
A2 1
This is exactly what we obtained in Section 12.2 for the firstpnormal mode of the
oscillating two-mass system. Similarly, the second eigenvalue ω2 = 3k/m, and the cor-
responding eigenvector
i i
i i
i i
A1 −1
=
A2 1
This is again the same result we obtained in Section 12.2 for the second normal mode.
2∗k −k
A= m
k
m
2∗k ;
−m m
Eigenvalues[A]//Simplify
OUTPUT: k
3k
m,m
Eigenvectors[A]//Simplify
1
A1
=
A2 1
OUTPUT
E i g e n v a l u e s= [ 3 ∗ k/m, k/m]
E i g e n v a l u e s= [ ( k/m, 1 , [ Matrix ( [ [ 1 ] , [ 1 ] ] ) ] ) ,
( 3 ∗ k/m, 1 , [ Matrix ( [ [ − 1 ] , [ 1 ] ] ) ] ) ]
and the Mathematica code produces the two eigenvalues and eigenvectors:
m1 −m2 +z
A1
s
k (m1 + m2 − z) m1
ω1 = =
m1 m2
A2 1
i i
i i
i i
Eigenvectors[A]//Simplify
√ √
m1−m2+ m12 −m1m2+m22 −m1+m2+ m12 −m1m2+m22
OUTPUT: m1 ,1 , − m1 ,1
1 1 1
T1 = m1 ẋ21 + ẏ12 = m1 L21 cos2 θ1 θ̇12 + L21 sin2 θ1 θ̇12 = m1 L21 θ̇12 (12.3.1)
2 2 2
The location (x2 , y2 ) of the second mass is shifted with respect to the first mass by (x1 , y1 ),
so that (x2 , y2 ) = (L1 sin θ1 + L2 sin θ2 , −L1 cos θ1 − L2 cos θ2 ). Therefore the kinetic energy
of the second mass is:
1
T2 = m2 ẋ22 + ẏ22 (12.3.2)
2
1
d 2 d 2
= m2 (L1 sin θ1 + L2 sin θ2 ) + (L1 cos θ1 + L2 cos θ2 ) (12.3.3)
2 dt dt
i i
i i
i i
+x
+y L1
θ1
m1
L2
θ2
m2
Figure 12.3: The double pendulum oscillator is characterized by the two angles (θ1 (t), θ2 (t)).
θ12
For amplitudes of small oscillations, we use the approximation cos θ1 ' 1 − 2 and cos θ2 '
θ22
1− 2 so that:
θ2 θ22
V = − (m1 + m2 ) gL1 1− 1 − m2 gL2 1 − (12.3.8)
2 2
The Lagrangian L = T − V of the system is:
1 1
L= (m1 + m2 ) L21 θ̇12 + m2 L1 L2 θ̇1 θ̇2 + m2 L22 θ̇22
2 2
θ12 θ2
+ (m1 + m2 ) gL1 1 − + m2 gL2 1 − 2 (12.3.9)
2 2
i i
i i
i i
−m2 L1 ω 2 A1 − m2 L2 ω 2 − g A2 = 0 (12.3.16)
This system of linear equations for A1 and A2 will have a nontrivial solution, only if the
determinant is zero:
(m1 + m2 ) −L1 ω 2 + g 2
−m 2 L2 ω
det =0 (12.3.17)
−m2 L1 ω 2 −m2 L2 ω 2 − g
In principle, we can obtain the solution to this characteristic equation, however the resulting
algebraic expressions are very complicated. So instead of looking at the completely general
solution, let us look at the special case of equal lengths (L1 = L2 = L), and identical masses
(m1 = m2 = m). In this special case, (12.3.17) becomes:
2m −Lω 2 + g 2
−mLω
det =0 (12.3.18)
−mLω 2 −m Lω 2 − g
i i
i i
i i
(a) (b)
Figure 12.4: The symmetric and antisymmetric normal modes of oscillation for the double
pendulum with equal masses and equal lengths. The two natural frequencies are ω1 =
q √ √
2 ± 2 g/L , and the amplitudes are related by A2 = ± 2A1 .
q √ √
The solution for ω1 = 2 − 2 g/L is then A2 = 2A1 , and we can write the first possible
solution corresponding to the first normal mode :
√ √
x1 (t) = A1 eiω1 t x2 (t) = 2x1 (t) = 2A1 eiω1 t (12.3.22)
√ √
Since A2 = 2A1 and θ2 (t) = 2θ1 (t), this type of motion corresponds to both masses
moving in the same direction and in phase at all times,√ as shown in Figure 12.4a. The
amplitude of the second pendulum must be equal to 2 times larger than the amplitude
of the first pendulum. This is the first normal mode or the symmetric mode of oscillation,
and the general motion of the system in this mode can be written as a linear combination
of trigonometric functions:
√
θ2 (t) = 2θ1 (t) cos (ω1 t − φ1 ) (12.3.23)
In matrix notation, the first normal mode can be written as:
θ1 (t) 1
= D1 √ cos (ω1 t − φ1 ) (12.3.24)
θ2 (t) 2
q √
By working in a similar manner with the second frequency ω2 = 2 + 2 g/L, we find:
θ1 (t)
−1
= D1 √ cos (ω1 t − φ1 )
θ2 (t) 2
This is the second normal mode or the antisymmetric mode of oscillation, in which both
masses moving in the same direction completely out of phase
√ at all times, as shown in Figure
12.4b. The amplitude of the second pendulum is again 2 times larger than the amplitude
of the first pendulum.
The above results can be obtained using the symbolic capabilities of Mathematica, by
using the Solve and Simplify commands. This is shown in Example 12.8.
i i
i i
i i
Example 12.8: The natural frequencies and normal modes of the double pen-
dulum
Find the natural frequencies and normal modes of the double pendulum with equal
masses and equal lengths.
Solution:
When m1 = m2 = m and L1 = L2 = L, the system of equations (12.3.15) and (12.3.16)
can be written in matrix form as:
−2m Lω 2 − g 2 0
−mLω A1
=
−mLω 2 −m Lω 2 − g 0
A2
Once more,qwe accept only the two positive values of the natural frequency ω. The solu-
√ √ √
tion for ω1 = 2 − 2 g/L has amplitudes related by A2 = 2A1 and θ2 (t) = 2θ1 (t).
q √ √
The solution for ω2 = 2 + 2 g/L has amplitudes related by A2 = − 2A1 and
√
θ2 (t) = − 2θ1 (t). These results are identical to the ones obtained using the analytical
approach.
2 ∗ m ∗ (−L ∗ ω ∧ 2 + g) −m ∗ L ∗ ω ∧ 2
A= ;
−m ∗ L ∗ ω ∧ 2 −m ∗ (L ∗ ω ∧ 2 − g)
A1 0
Solve A. == , {ω, A1, A2} //Simplify
A2 0
OUTPUT:
√
(−2+ 2)g √
q
{A1 → 0, A2 → 0}, ω → − − L , A2 → 2A1 ,
√
(−2+ 2)g √
q
ω→ − L , A2 → 2A1 ,
n p √ qg √ o n p √ qg √ oo
ω → − 2+ 2 L , A2 → − 2A1 , ω → 2 + 2 L , A2 → − 2A1
By substituting θ̈1 from this expression in (12.3.11), we obtain the equation for θ̈2 :
i i
i i
i i
(m1 + m2 ) g m2 g
− A1 ω 2 = − A1 + A2 (12.3.27)
m1 L1 m1 L1
(m1 + m2 ) g g (m1 + m2 )
− A2 ω 2 = − A1 − A2 (12.3.28)
m1 L2 m1 L2
In matrix form:
(m1 +m2 )g
− mm12Lg1
m1 L 1 A1 A1
= ω2 (12.3.29)
(m1 +m2 )g
m1 L 2 − g(m 1 +m2 )
m1 L 2
A2 A2
This is once more the standard form of the eigenvalue problem, and we can proceed by
finding the eigenvalues and eigenvectors of the 2 × 2 matrix on the left hand side of this
matrix equation. Finding the natural frequencies and the amplitudes of oscillation A1 and
A2 in Python and Mathematica is left as an exercise, see the problems at the end of this
chapter.
rα = rα (q1 , . . . , qn ) , (12.4.2)
We then showed that the kinetic energy which is written in Cartesian coordinates as:
1X
T= mα ṙα · ṙα (12.4.3)
2 α
i i
i i
i i
where
X ∂rα ∂rα
ajk = mα · (12.4.5)
α
∂qj ∂qk
and the coefficients ajk = ajk (q1 , . . . , qs ) depend on the coordinates qj .
Let us now assume that the potential energy of the system depends only on the coordi-
nates qj :
V = V (q1 , . . . , qs ) (12.4.6)
Under the above assumptions, our Lagrangian has the form:
1X
L = T −V = ajk q̇j q̇k − V (q1 , . . . , qs ) (12.4.7)
2
j,k
Let us now assume that the system performs small oscillations around a point of stable
equilibrium qEQ = (q1 , q2 , . . . , qs )EQ . Without loss of generality, we can redefine our coor-
dinates so that this equilibrium point is the origin of our generalized coordinate system,
i.e., qEQ = 0. Next, we expand the potential energy around this point:
s s
X ∂V 1 X ∂2V
V = V (0) + qi + qj qk + ... (12.4.8)
∂qi 2 ∂qj ∂qk
i=1 j,k=1
where all derivatives are evaluated at the equilibrium point located at the origin of our
coordinate system, qEQ = 0. At the equilibrium points we have ∂V /∂qi = 0, and we can
further simplify this expression by setting V (0) = 0 to obtain:
s
1 X
V = Kjk qj qk (12.4.9)
2
j,k=1
s
1 X
V = Kjk qj qk Kjk = constants (12.4.12)
2
j,k=1
i i
i i
i i
s s
1 X 1 X
L = T −V = ajk q̇j q̇k − Kjk qj qk (12.4.13)
2 2
j,k=1 j,k=1
Example 12.9 shows how to calculate the coefficients ajk and Kjk for the double pendulum
system we studied in the previous section.
Solution:
The kinetic energy of the double pendulum for small oscillations was found in the
previous section to be:
1 1
T= (m1 + m2 ) L21 θ̇12 + m2 L1 L2 θ̇1 θ̇2 + m2 L22 θ̇22
2 2
1 1 1 1
= (m1 + m2 ) L1 θ̇1 + m2 L1 L2 θ̇1 θ̇2 + m2 L1 L2 θ̇1 θ̇2 + m2 L22 θ̇22
2 2
2 2 2 2
The coefficients ajk are the coefficients appearing in (12.4.4) for the kinetic energy
T . Since our generalized coordinates are q1 = θ1 and q2 = θ2 , then ajk are the coeffi-
cients of the θ̇12 /2, θ̇22 /2, and θ̇1 θ̇2 /2 terms in this expression for T . Therefore a11 =
coefficient of θ̇12 /2 = (m1 + m2 ) L21 , a12 = a21 = coefficient of θ̇1 θ̇2 /2 = m2 L1 L2 , and
The potential energy of the double pendulum was found in the previous section to be:
θ12 θ22
V = − (m1 + m2 ) gL1 1 − − m2 gL2 1 −
2 2
The coefficients Kjk are the coefficients appearing in (12.4.9) for the potential energy
V , for the θ12 , θ22 , and θ1 θ2 terms. Therefore K11 = coefficient of θ12 = (m1 + m2 )gL1 ,
K12 = K21 = coefficient of (θ1 θ2 ) = 0, and K22 = coefficient of θ22 = m2 gL2 .
We can then use the constant coefficients ajk and Kjk as two 2 × 2 matrices. The
coefficients ajk are used to create the matrix M and the coefficients Kjk are used to
create the matrix K:
As we will see next, the matrix M can be thought of as a generalized “mass matrix”
and the matrix K as a generalized “spring constant” matrix K.
12.4.2 The Equations of Motion for Small Oscillations Around an Equilibrium Point
We can now proceed to derive the general equations of motion, by evaluating the Euler-
Lagrange expressions:
s
∂L ∂V X
=− =− Kjk qk (12.4.14)
∂qj ∂qj
k=1
i i
i i
i i
s n
!
d ∂L d ∂T d X X
= = ajk q̇k = ajk q̈k (12.4.15)
dt ∂ q˙j dt ∂ q˙j dt
k=1 k=1
Finally, we can write these equations of motion in compact form as a single matrix equation.
a11 a12 ··· a1s q̈1 K11 K12 ··· K1s q1
a21 a22 ··· a2s q̈2 K21 K22 ··· K2s q2
.. .. .. .. .. = .. .. .. .. .. (12.4.18)
. . . . . . . . . .
as1 as2 ··· ass q̈s Ks1 Ks2 ··· Kss qs
with the s × 1 column matrix q denoting the generalized coordinates of the system, M rep-
resenting the symmetric s × s generalized “mass matrix”, and K the symmetric generalized
“spring constant” matrix.
q1 a11 . . . a1s K11 . . . K1s
q = ... M = ... ..
.
.. K = ..
. . ..
.
..
. (12.4.19)
Since we expect oscillatory motion, we try solutions of the form qj (t) = Aj eiωt , where Aj are
the amplitudes of oscillation and ω is the natural frequency of the system. By substituting
into (12.4.17) and canceling the exponential terms as usual, we obtain:
ω 2 MA = KA (12.4.20)
with the s × 1 column matrix A (whose elements are the Aj from the trial solution qj )
denoting the unknown amplitude matrix.
Equation (12.4.20) is the generalized matrix form of the familiar equation mω 2 = k,
which describes the oscillation frequency ω for a mass m attached to a spring with spring
constant k. This equation can be solved in principle by using any of the linear algebra
techniques we saw previously in this book.
Example 12.10 shows one method of solving (12.4.20) for the double pendulum.
Example 12.10: Finding the natural frequencies and normal modes of the dou-
ble pendulum, again
Solve ( (12.4.20) for the double pendulum with equal masses and equal lengths, and
find the natural frequencies and the normal modes.
i i
i i
i i
Solution:
In Example 12.9, we calculated the coefficients ajk and Kjk of the matrices M and K
for the double pendulum. Equation (12.4.20) then becomes:
2mL Lω 2 − g mL2 ω 2 0
A1
=
mL2 ω 2 mL Lω 2 − g A2 0
This, of course, is the same equation we obtained and solved in Example 12.8 (to see
this, divide the above by −L), and we can reuse the code for Example 12.8 to find ω and
A.
In the next subsection, we proceed to show how the above general matrix formalism of the
equations of motion leads to the concept of normal coordinates.
i i
i i
i i
Since the vectors are linearly independent, the last equation can hold only if the coeffi-
cients of the vector MAj are equal:
In words, this equation tells that any solution q of the equations of motion, can be
written as the linear combination of s independent oscillating terms, each of which has its
own characteristic natural frequency ωj . These oscillating terms are known as the normal
coordinates cj (t) of the system. Hence, all motions of a linear coupled oscillator can be
expressed as a linear combination of normal coordinates.
The material in this section is a formal way of demonstrating the result of Example 12.4.
In Example 12.4 we found that the analytical solution, of the case
equal masses
p andequal
springs, could be written as the sum of two terms, cos k/mt and cos 3k/mt . For
p
example, the
p solution x1 (t) could be written in the form (12.4.29) by writing ω1 = k/m
p
0
a/2
A1 = A2 = (12.4.30)
0 −a/2
as the amplitudes.
i i
i i
i i
x1 (t) 1
= D1 cos (ω1 t − φ1 )
x2 (t) 1
The second possible p solution for the motion of the two masses corresponds to the
natural frequency ω2 = 3k/m, and in this type of motion the two masses moves in
opposite directions, while the center of mass remains stationary.
This type of motion is known as the second normal mode or the antisymmetric mode
of oscillation, and can be written as:
x1 (t) 1
= E1 cos (ω2 t − φ2 )
x2 (t) −1
In general, the motion of the system will be a linear p combination ofpthe two possible
normal modes, corresponding to the frequencies ω1 = k/m and ω2 = 3k/m.
This system can also be reduced into an eigenvalue problem, by substituting a trial
solution of the form x1 (t) = A1 eiωt and x2 (t) = A2 eiωt in the system of coupled equations.
In a second well-known example of coupled harmonic oscillators, we examined the
double pendulum, which also has two normal modes. The frequencies ofqthe symmetric
q √ g √ g
and antisymmetric modes of oscillation are ω1 = 2− 2 L and ω2 = 2+ 2 L .
In this chapter, we also looked at the general mathematical formalism for the normal
modes of coupled oscillators. Here we use generalized coordinates qj (j = 1, 2, . . . , s) and
the corresponding generalized velocities q̇j , where s is the number of degrees of freedom
of the system.
The Lagrangian for small oscillations around an equilibrium position can be reduced
into a quadratic function of the generalized coordinates and generalized velocities of the
system in the form:
n
1X 1 X
L = T −V = ajk q̇j q̇k − Kjk qj qk
2 2
j,k j,k=1
i i
i i
i i
This system of coupled differential equations can be solved using standard techniques
of linear algebra.
Any solution q = qj (j = 1, 2, . . . , s) of these equations of motion can be written as the
linear combination of s independent oscillating terms, each of which has its own charac-
teristic natural frequency ωj :
s
X
q= bj sin (ωj t − φj ) Aj
j=1
These oscillating terms are known as the normal coordinates of the system.
2. Consider the two-mass three-spring system studied in this chapter, with different
masses m1 , and m2 and identical spring constants k1 = k2 = k3 .
(a) Transform the Lagrangian by introducing two new variables q1 = x1 +x2 and q2 =
x1 − x2 . Obtain new equations of motion by using the Euler-Lagrange equations
for this transformed Lagrangian.
(b) Solve the new equations and obtain q1 (t) and q2 (t), and the natural frequencies
of the system. This is an example where by using a new set of coordinates, we
can uncouple the differential equations of motion. These new coordinates q1 (t),
q2 (t) represent the normal coordinates of the system.
3. In the two-mass three-spring system studied in this chapter, obtain the analytical
solution x1 (t) and x2 (t), by using the following initial conditions: the first mass is
at rest in the equilibrium position, and the second mass is moved a distance a from
equilibrium, and released from rest.
4. Obtain an approximate analytical solution for the case of two weakly coupled oscilla-
tors described in Example 12.5. This system is weakly coupled if the spring constant
k2 of the middle spring is much smaller than the spring constants of the two end
springs (k2 << k1 = k3 ). Use m1 = m2 and the following initial conditions: the first
mass is at rest in the equilibrium position, and the second mass is moved a distance
a from equilibrium, and released from rest. Find the frequency ω at which energy is
transferred back and forth between the two coupled oscillators.
5. Consider two equal masses m attached to one wall and to each other with springs of
spring constant √
k, as shown in Figure 12.5. Show that the natural frequencies of the
system are ω = 5±1
p
2 k/m.
i i
i i
i i
6. Consider the three different situations shown in Figure 12.6, where three identical
masses m are attached to springs with the same spring constant k. Write the equations
of motion for these three cases.
7. Three equal masses m are connected to each other with identical springs of constant
k, as shown in Figure 12.7. Find and describe the normal modes of oscillation (assume
no friction), by using appropriate programming code.
Figure 12.7: Problem 12.7: Three identical masses connected by two identical springs.
8. Three equal masses m are connected with springs of the same spring constant k.
The masses are connected to each other and to two end walls, as shown in Figure
12.8. Calculate the frequencies of the normal-modes of oscillations of this system and
describe the corresponding motion of the three masses in each normal mode.
i i
i i
i i
Figure 12.8: Problem 12.8: Three identical masses connected to four identical springs.
9. Find the effective spring constants and natural frequencies for each of the oscillating
systems shown in Figure 12.9 below. In Figure 12.9, two springs with spring constants
k1 and k2 are attached in parallel to a mass m and to one wall. In Figure 12.9b, two
springs with spring constants k1 and k2 are attached in series to a mass m and to one
wall.
(a) (b)
10. Two identical masses are connected with identical springs with spring constant k, and
are restricted to move on a circle as shown in Figure 12.10. Find expressions for the
natural frequencies and describe the motion of the system.
11. Repeat Problem 10, with three identical masses connected with identical springs with
spring constant k.
i i
i i
i i
12. Four identical masses are connected by four identical springs, and are constrained to
move on a frictionless circle of radius R. How many normal modes of small oscillations
are there and what are the frequencies of small oscillations?
13. Three bodies of equal mass m are connected with springs of constant k, and are placed
in an equilateral shape as shown in Figure 12.11. The masses are constrained to move
on the xy-plane. Find the natural frequencies of this system.
14. A simple classical model of the CO2 molecule would be a linear structure of three
masses, with the electrical forces between the ions represented by two identical springs
of equilibrium length L and force constant k. Assume that only motion along the
original equilibrium line is possible, i.e., ignore rotations. Let m be the mass of each
oxygen and M be the mass of the carbon atom
(a) How many vibrational degrees of freedom does this system have?
(b) Find the normal modes and calculate the natural frequencies.
15. Two equal masses m are suspended from a ceiling by springs with constants k, as
shown in Figure 12.12.
(c) Find the equations of motion of the two masses by using a coordinate system
(y1 , y2 ) as measured from the ceiling.
(d) Transform these equations of motion by making a substitution y1 = x1 + u1 and
y2 = x2 + u2 , where u1 , and u2 are the equilibrium positions of the two masses.
Solve the new system of equations and described the normal modes of oscillations
of this system in the vertical direction.
Figure 12.12: Problem 12.15: Two masses hanging by springs from a vertical support.
i i
i i
i i
Figure 12.13: Problem 12.17: Two unequal mass pendulums connected by a spring.
i i
i i
i i
19. A mass M moves along the x-axis on a frictionless surface, and is attached to a
wall with a spring with spring constant k as shown in Figure 12.15. The mass m is
connected to M by a massless string, as shown below.
(a) Write a Lagrangian for this system for a small angle approximation.
(b) Find x(t) and θ(t)
(c) Find the frequencies of small oscillations.
20. A thin uniform bar of mass m and length 3L/2 is suspended by a string of length L
and negligible mass as shown in Figure 12.16. Find the normal frequencies for small
oscillations in a plane.
21. A mass m moves in a gravitational field g pointing in the z-direction, and on the
inside wall of a frictionless axially symmetric vessel given by z = b(x2 + y 2 )/2, where b
is a constant and z is in the vertical direction, as shown in Figure 12.17. The particle
is moving in a circular orbit at height z = zo .
(a) Obtain the angular frequency of the circular motion in terms of m, zo , b, and g
(b) Find the total energy and angular momentum in terms of m, zo , b, and g
(c) The particle is pushed slightly downwards in the horizontal circular orbit. Obtain
the frequency of oscillation about the unperturbed orbit, in the case of small
oscillations around the stable circular motion.
i i
i i
i i
Figure 12.17: Problem 12.21: An axially symmetric vessel given by z = b(x2 + y 2 )/2.
23. A hemisphere with a nonuniform density is placed on a table as shown in Figure 12.18.
The radius of the hemisphere is R, and the center of mass is located a distance a above
the table.
(a) Find the potential energy when the hemisphere is slightly misplaced by an angle
θ from equilibrium. Under what conditions will the object perform oscillations
around a stable equilibrium point?
(b) Find the frequency of small oscillations around this stable equilibrium point,
when the moment of inertia of the object around the center of mass is ICM .
Assume that the object is only allowed to roll around the pivot contact point.
i i
i i
i i
24. A uniform sphere of radius r rests in equilibrium inside a uniform hemispherical shell
of radius R, as shown in Figure 12.19.
(a) Find the relationship between r and R in order to have a stable equilibrium
situation for small oscillations.
(b) Find the natural frequency for small oscillations of the system.
i i
i i
i i
CHAPTER 13
Nonlinear Systems
Nonlinear systems are everywhere! In physics, one of the most common examples of a
nonlinear system is an oscillator undergoing large amplitude oscillations. However, nonlin-
ear systems occur in other sciences as well. Examples of nonlinear systems in other fields
include predator-prey models, population dynamics, disease propagation, arms-race models,
economic models, and the list goes on. In fact, nonlinear systems are the norm and linear
systems are often approximations. While the diversity and difficulty of nonlinear systems
may be concerning, there are some basic methods of analysis which are applicable to most
of the nonlinear systems you would encounter. In this chapter, we will learn some of those
methods. Along the way, we will discover some interesting new behaviors that only nonlin-
ear systems can exhibit such as bifurcations and chaos. At the end of the chapter, you will
have an opportunity to apply what you have learned to analyze systems from a variety of
natural and social sciences.
407
i i
i i
i i
still considered a linear system because all of the terms are linear in the dependent variable
x and its derivatives.
One of the first consequences of nonlinearity is that the principle of linear superposition
no longer holds. Consider the equation of motion for a simple harmonic oscillator, ẍ + ω02 x =
0, studied in Chapter 6. Suppose we find two solutions for this equation, x1 (t) and x2 (t).
Then, the general solution to the harmonic oscillator equation is: c1 x1 (t) + c2 x2 (t), where
c1 and c2 are constants. We can see that by inserting the sum into the differential equation
ẍ + ω02 x = 0:
Now, let us change the differential equation to the form: ẍ + x2 = 0, and say we find two
solutions to this equation that we will also call x1 and x2 . Is it true that c1 x1 (t) + c2 x2 (t)
is still a solution? We insert the superposition into the nonlinear differential equation:
(c1 x¨1 + c2 ẍ2 ) + (c1 x1 + c2 x2 )2 = c1 x¨1 + x21 + c2 ẍ2 + x22 + 2c1 c2 x1 x2 = 2c1 c2 x1 x2
Notice that now, we do not get zero. In other words, the differential equation no longer
satisfies superposition. The result is that unlike linear systems, nonlinear systems cannot be
broken into parts, solved individually, and then have those solutions combined for a solution
to the whole system. Recall the topic of the driven damped harmonic oscillator. If there are
multiple sinusoidal drive terms on the righthand side of the equation, we could solve the
differential equation for one right-hand term at a time, then add the individual solutions to
get a particular solution. That cannot be done for nonlinear systems. This makes nonlinear
systems more difficult to solve in closed form. In fact, closed form solutions for nonlinear
differential equations often cannot be found.
At this point, we might be wondering if any of this matters. Are nonlinear systems
just mathematical curiosities? The answer is no. In fact, many systems from the physical,
biological, and social sciences are nonlinear, and linear descriptions of behavior are often
approximations. For example, Hooke’s law, F (x) = −kx, holds for an oscillator undergoing
small amplitude oscillations. Once the amplitude of the oscillations is large enough, then
the restoring force is no longer linear, and nonlinear terms need to be included in order for
the equation of motion to accurately describe the oscillation. One of the results is that the
period of oscillation is no longer independent of the amplitude of oscillation.
As we will see, nonlinear systems can display a much more diverse set of behaviors than
linear systems, which tend towards one of four basic types of behavior: exponential growth,
exponential decay, oscillation, and a combination of the previous three. The limited types of
behavior and mathematical simplicity makes linear systems easier to work with and allows
us to predict future states of the system with a high degree of accuracy at any point in
the future. Nonlinear systems, however, are often not exactly solvable and display complex
behaviors such as bifurcations and chaos, which can make long term predictions of behavior
impossible.
So, what do we do if we encounter a nonlinear system? Do we attempt to make a linear
approximation and hope that the approximation is good enough? Do we just give up because
we cannot solve the problem? First, we do not have to give up because we cannot solve a
problem in closed form. The numerical methods described in this book so far are powerful
tools and can provide insight into the behavior of a nonlinear system. In this chapter, we
will expand upon those numerical methods, learning how to interpret these solutions in new
ways. Second, linear approximations can work well and can shed some light on the solution
i i
i i
i i
of a nonlinear system. These linear approximations, often done using Taylor expansion, will
allow us to use some of the tools for linear differential equations in order to provide some
understanding of the nonlinear system. Ultimately, one of the most powerful tools we have
in our arsenal is the phase space plot mentioned in Chapter 6. The phase space plots will
provide graphical methods of understanding the long-term behavior of a nonlinear system.
Recall, that for the damped driven harmonic oscillator, we were mostly interested in the
steady-state solution and less interested in the transient behavior. The same will generally
be true for nonlinear systems, and phase space plots will allow us to identify the nature of
the steady-state solution.
Beyond physics, nonlinear terms appear in equations that govern systems which have
interacting components. For example, equations that govern population size, predator-prey
interactions, certain chemical reactions, arms races, and the spreading of disease (just to
name a few) are nonlinear. Therefore, the mathematical techniques we will learn in this
chapter are applicable to systems from a wide range of natural and social sciences. Some
of the most exciting developments in science is the application of “traditional physics tech-
niques” to problems in the biological and social sciences.
i i
i i
i i
import numpy a s np
from s c i p y . i n t e g r a t e import o d e i n t
import m a t p l o t l i b . p y p l o t a s p l t
d e f dho ( y , t ) : #d e f i n e t h e ODEs t o be s o l v e d
x, v = y
gamma = 0 . 1
d e r i v = [ v , −2.0∗gamma∗v − x ] #RHS o f ODEs
return deriv
x v e c = [ 1 . 0 , 0 . 0 ] #x ( 0 ) = 1 and #v ( 0 ) = 0
t = np . l i n s p a c e ( 0 , 1 0 0 , 1 0 0 1 ) #time r a n g e f o r t h e s o l u t i o n
s o l u t i o n = o d e i n t ( dho , x ve c , t )
plt . plot ( solution [ : , 0 ] , solution [ : , 1 ] )
Algorithm 12: The Python code used to generate the phase space diagram for (13.2.2).
The results of this code for γ = 0.1 and γ = 0 are shown in Figure 13.1. Pretend, for a
moment, that we did not know how to solve this problem in closed form, but instead all
we had were the phase space plots that were generated numerically; this is often true for
nonlinear systems. What can we learn from Figure 13.1 about the behavior of (13.2.1)? Let
us first consider the left graph in Figure 13.1, where γ = 0.1 and, hence, there is a small
amount of damping in the system. To help visualize the behavior of the system, we will
think of (13.2.1) as describing a mass on a spring with damping. The path in phase space,
called the phase space trajectory, or simply trajectory, describes the behavior of the system,
and represents the solution to (13.2.2) and therefore, also to (13.2.1). To understand the
behavior of the system, follow the trajectory in phase space! As we follow the trajectory from
its starting point we see that the position decreases towards zero while the velocity becomes
i i
i i
i i
γ = 0.1 γ=0
1.00
1.0
0.75
0.50
0.5
0.25
velocity
velocity
0.00 0.0
−0.25
−0.5
−0.50
−0.75
−1.0
−1.00
−1.0 −0.5 0.0 0.5 1.0 −1 0 1
position position
Figure 13.1: The phase space plot for (13.2.2) with γ = 0.1 (left) and γ = 0 (right).
more negative. Thinking of the a mass on a spring, we see the mass starting at rest at a
position of 1.0 m to the right of equilibrium (x = 0), and when released, it moves to the left
towards equilibrium, gaining speed. The mass reaches equilibrium where it has its maximum
(negative) speed, overshoots equilibrium, and begins to slow down (velocity becomes less
negative). The mass continues moving to the left to a distance of about 0.6 meters left
of equilibrium and, at that point, the velocity is zero. The mass, however, has reached a
turning point and moves to the right (positive velocity) gaining speed as it approaches
equilibrium again. The mass reaches its highest speed at equilibrium and beings to slow
down. However, notice that the mass does not return to its initial position, instead ending
up at a distance of about 0.5 m to the right of equilibrium before turning around. This
decaying oscillation continues until the mass is at rest at equilibrium (trajectory ends at
the origin).
Hence, we can use Figure 13.1 to answer our question about the steady-state of the
system (the mass is at rest at the equilibrium position) and we could even discuss the
transient behavior, all without having a closed form solution! The final point, in this case
the origin, is a point where both ẋ = 0 and v̇ = 0. Points in phase space that satisfy the
conditions, ẋ = 0 and v̇ = 0, are called fixed points of the system. Fixed points correspond to
equilibrium points of the system and take on different forms. In the case of γ = 0.1, the fixed
point is called a stable spiral and it represents a decaying oscillation onto the equilibrium
position. There are many different types of fixed points and we will discuss how to identify
and classify some of them. However, a thorough discussion of the techniques for finding
and classifying fixed points are beyond the scope of this book. For such a discussion, the
interested reader is directed to [Strogatz(2014)].
We can get the same qualitative information about the system’s behavior in the case of
γ = 0,which corresponds to the right graph in Figure 13.1. Again, starting with an initial
condition of (x = 1, v = 0), and continuing with the mass-spring analogy, we can follow the
trajectory clockwise to see that the mass moves towards equilibrium with an increasing
speed (to the left), overshoots equilibrium, slows down and stops a distance of 1.0 m to
the left of equilibrium. At that point, the mass begins moving to the right with increasing
speed until it overshoots equilibrium and returns to its starting position. The closed curve
trajectory represents an oscillation, the system returning to its initial state. The origin is
i i
i i
i i
still an equilibrium position and is still a fixed point, but the mass does not come to rest
there. In this case, the fixed point at the origin is referred to as a center (as in center
of oscillation). If the initial condition was the origin, the mass would stay there forever,
without moving. This is the same as placing the mass at rest at the origin.
As we can see, it is possible to get a lot of valuable information from phase space plots.
We can ask further questions. For example, do all initial conditions have trajectories that
spiral into the origin when γ = 0.1? The answer, of course, is yes because damping removes
energy from the system and causes it to come to rest a equilibrium. However, we could find
this out by selecting several initial conditions and numerically solving for those using the
code in Algorithm 12. We would find that all points have trajectories that spiral into the
origin and therefore the origin is a global attractor, a fixed point which “attracts” trajectories
from all points in the phase space. Likewise, we could repeat the exercise for γ = 0 and find
that the resulting phase space graph consists of nested closed trajectories. In this case, the
origin neither attracts nor repels trajectories and is considered to be neutrally stable. Hence,
we could answer the question about the behavior of the system for any initial condition,
even without solving the system in closed form.
In the above example, we saw two types of fixed points, the stable spiral and a center,
representing a stable and neutral equilibrium, respectively. There are other types of fixed
points which repel trajectories, corresponding to unstable equilibrium, as well as fixed points
that both attract and repel trajectories (also corresponding to unstable equilibrium). In
the next section, we will discuss fixed points more generally. We will look into how to
find them and how to classify the stability of their equilibrium states. The next section is
mathematical in nature and a bit abstract, but we will follow with two examples which will
help demonstrate the key ideas of the section.
where f (x, y) and g(x, y) are functions of the variables x and y. Recall that in (13.2.2)
we used the variables x and v instead of x and y in order to reinforce the idea that our
variables represent position and velocity. Using generic variables x and y, (13.3.1) allows
for the analysis of more general systems. For example, x and y could represent position and
velocity for a physical system, or they could represent the numbers of rabbits and foxes in
a predator-prey model. We can recast (13.2.2) into a form similar to (13.3.1) by choosing
y = v and (13.2.2) becomes,
ẋ =y
ẏ = − 2γy − x
where f (x, y) = y and g(x, y) = −2γy − x. Equations (13.3.1) are sometimes referred to as
a two-dimensional system of equations because they depend on two variables, x and y. In
i i
i i
i i
addition, (13.3.1) is called the equations of motion for our system. Recall that the term,
equations of motion, for a physical system describes the actual motion of a particle. However,
it is common to refer to (13.3.1) as equations of motion even if we are using them to describe
the populations of animals, or any other type of system that is not necessarily moving (but
is changing). To put it more simply, the equations of motion describe the behavior of a
given system regardless of the nature of the system.
While nonlinear systems are diverse and often difficult to analyze, there are some general
strategies on how to approach such systems. For example, a general approach that we like
to recommend to our students is the following:
The General Approach to Analyzing Nonlinear Systems
1. Find the fixed points (equilibrium points) of the system, points (x∗ , y ∗ ) such that
f (x∗ , y ∗ ) = g (x∗ , y ∗ ) = 0.
2. Identify the stability of the fixed points.
3. Using initial conditions near the fixed points, plot the phase portrait, a plot of the
xy-plane in phase space, which contains a representative selection of trajectories.
4. Identify other nonlinear structures in the phase portrait such as limit cycles and
strange attractors (see below).
5. Identify bifurcations (see below).
In this section, we will focus on the first two steps, and then, in later sections, we will
demonstrate the final two steps by example.
To find the fixed points, (x∗ , y ∗ ) for the system (13.3.1), we simply solve the equations:
)
f (x∗ , y ∗ ) =0
(13.3.2)
g (x∗ , y ∗ ) =0
for (x∗ , y ∗ ). In our damped harmonic oscillator example, there is one fixed point
(x∗ = 0, y ∗ = 0), found by solving f (x∗ , y ∗ ) = y ∗ = 0 and g (x∗ , y ∗ ) = −2γy ∗ − x∗ = 0. Sys-
tems may have multiple fixed points. Solving for fixed points can be done either by hand,
or by using a computer algebra system such as Mathematica’s Solve command. Recall that
these fixed points correspond to equilibrium states of the system, for example, the mass
being at rest at the origin.
Once we have found the fixed points, we need to find the stability of those fixed points. In
other words, are the equilibrium states stable, unstable, or neutrally stable? It turns out that
there is a method of calculating the stability of the equilibria for two-dimensional systems.
To find the stability, we need to know how the system behaves near each equilibrium state,
i.e., we need to find the mathematical form of the equations of motion near the equilibrium
point. Taylor series expansions provide a polynomial approximation of equations near a
particular point. Hence, by Taylor expanding f and g, we can get a simple mathematical
form for the equations of motion near the equilibrium states. The Taylor-expanded f and
g will result in simple linear differential equations that can be solved in closed form. The
closed form solutions will tell us the stability of the system near the equilibrium point. In
summary, in order to get the stability of the fixed point, Taylor expand the equations of
motion about the fixed point and solve the resulting differential equations.
To perform the Taylor series expansion, we perturb the system about the fixed point.
The perturbation is done via a change of coordinates: x = x∗ + u and y = y ∗ + v, where u
i i
i i
i i
and v are small compared to x∗ and y ∗ , respectively. Note that v is not being used to denote
velocity. Furthermore, the origin of our new coordinates, (u, v), is the fixed point (x∗ , y ∗ ).
Hence, when we say we have perturbed the system about the fixed point, we can think of
starting the system just slightly off of the fixed point (equilibrium position) in phase space.
Next, we substitute x = x∗ + u and y = y ∗ + v into the equations of motion (13.3.1) and
Taylor expand about the point (x∗ , y ∗ ). Note that ẋ = u̇ because x∗ is constant, similarly
for ẏ.
After substitution, the Taylor expansion of f in (13.3.1) becomes:
u̇ =f (x∗ + u, y ∗ + v) (13.3.3)
∂f ∂f
=f (x∗ , y ∗ ) + u +v +O u2 , v 2 (13.3.4)
∂x (x∗ ,y∗) ∂y (x∗ ,y∗)
∂f ∂f
=u +v (13.3.5)
∂x (x∗ ,y∗) ∂y (x∗ ,y∗)
Notice that some terms appear to be lost from (13.3.4)to (13.3.5). The term f (x∗ , y ∗ ) =
0 by the definition of fixed points, while the term O u2 , v 2 contains second-order and higher
terms. Recall that u and v are small, so u2 and v 2 are very small and negligible for the
dynamics. While there are cases where those nonlinear terms need to be retained, we will
focus this text on systems for which those terms can be neglected. The same procedure can
be repeated for g(x, y):
∂g ∂g
v̇ = u +v (13.3.6)
∂x ∗ (x ,y∗)∂y ∗ (x ,y∗)
Hence, we find that by Taylor expanding (13.3.1) about the fixed point, we obtain a
system of equations that is linear in u and v. The above process is sometimes called the
linearization of the equations of motion. We can solve these linearized equations in closed
form and get an understanding of how the trajectories behave near the fixed point. In other
words, we will be able to find the stability of the equilibrium point. It is helpful to cast the
linearized equations into matrix form:
where we have removed the notation for the evaluation of each partial derivative for sim-
plicity in notation. The 2 × 2 matrix in (13.3.7) is called the Jacobian matrix A and, for
simplicity, we can recast (13.3.7) in the form,
u̇ = Au (13.3.8)
where u = (u, v). Equation (13.3.8) represents a system of first-order differential equations
which can be solved with the solution u = weλt , where w is a constant vector. Inserting our
trial solution into (13.3.8) yields,
λw =Aw (13.3.9)
i i
i i
i i
which is the eigenvector equation for the matrix A. Hence, in our trial solution, u = weλt , λ
is the eigenvalue of A associated with the eigenvector, w. We have previously studied eigen-
values and eigenvectors in Chapters 11 and 12. Here we will study them from a geometric
perspective.
Returning to our analysis, we can rewrite (13.3.9) in the form,
(A − λ1) w = 0 (13.3.10)
where 1 is the 2 × 2 identity matrix (with 1’s on the diagonal and 0’s otherwise) and 0 is the
zero vector. The field of Linear Algebra tells us that if the matrix, A − λ1, is invertible, then
the solution to (13.3.10) is w = 0, which is the trivial solution. For a nontrivial (and therefore
interesting!) solution, we need the matrix A − λ1 to be singular (i.e., not invertible). Linear
Algebra also tells us that a matrix is singular if its determinant is zero. Hence, in order
to get a nontrivial solution, det (A − λ1) = 0 must be true. The condition det (A − λ1) = 0
produces a polynomial in λ, called the characteristic polynomial, which can be solved to
obtain the eigenvalues of A. Because the matrix, A, is a 2 × 2 matrix, it will produce
two eigenvalues and each eigenvalue can be plugged into (13.3.10) to get its associated
eigenvector. Using the principle of linear superposition, the solution to (13.3.8) is therefore,
i i
i i
i i
case represents an unstable equilibrium since perturbations from the equilibrium (the fixed
point) grow exponentially as t → ∞.
At this point, the math might be overwhelming. The procedure is not difficult (we
promise!) and can be performed to solve problems even if the linear algebra isn’t fully
understood (although you should study linear algebra, it is really quite interesting). We
will illustrate this procedure using example problems in order to clarify the steps outlined
above.
Example 13.1: The damped harmonic oscillator, revisited, again
Consider the equations of motion for the damped harmonic oscillator with m = 1:
ẋ =y
ẏ = − 2γy − ω02 x
Find the Jacobian matrix, compute its eigenvalues and its eigenvectors.
Solution:
We have already noted that f = y and g = −2γy − x and that the fixed point is located
at the origin. We can use (13.3.7) to get the Jacobian,
∂f ∂f
0 1
∂x ∂y
A= =
∂g ∂g
−ω 2 −2γ
∂x ∂y (0,0) 0
0 1 1 0 1
−λ
det(A − λ1) = det = =0
−λ
−ω02 −2γ 0 1 −ω02 −2γ − λ
0 1 1 0 0
w1
q
(A − λ1) w = − −γ + γ − ω0
2 2 =
−ω02 −2γ 0 1 w2 0
or
i i
i i
i i
q
γ − γ 2 − ω02 w1 + w2 =0
q
−ω02 w1 + −γ − γ 2 − ω02 w2 =0
resulting in
q
−1 γ + γ 2 − ω02
w1
= ω02
w2
1
What does all of this say about the stability of the equilibrium at the origin? We will
return to that later in this section, for now, this example is intended to simply demonstrate
the procedure outlined so far in this section.
This next example is not physically motivated. In other words, the equations do not
describe a particular physical system. However, it is a simple system of equations that allows
for a direct demonstration of the procedure and is more easily extended to a discussion of
stability than the example above.
ẋ =y(x + 1)
ẏ =x(y + 1)
Find the fixed points and solve the corresponding linearized system of equations.
Solution:
First, we need to solve for y(x + 1) = 0 and x(y + 1) = 0. It should be easy to verify for
this problem that there are two fixed points, (x∗ = 0, y ∗ = 0) and (x∗ = −1, y ∗ = −1). For
this example problem, we will study the fixed point at the origin and leave the analysis of
the other fixed point for the reader.
Following the procedure outlined above, we would perturb our solution using, x∗ = 0+v
and y ∗ = 0 + v, plug into our system of equations and Taylor expand to get the linearized
system. Fortunately, however, we have already done all of that work above and can simply
plug into (13.3.7):
x+1
u̇ y u
=
v̇ y+1 x v
(x=0,y=0)
0 1
u̇ u
=
v̇ 1 0 v
Now that we have our Jacobian, we can compute the eigenvalues by finding the char-
acteristic polynomial using, det(A − λ1) = 0:
0−λ 1
1 = λ2 − 1 = 0
0−λ
i i
i i
i i
0 1 1 0 −1 1
w1 w1
− = =0
1 0 0 1 w2 1 −1 w2
The above equation corresponds to the system of equations:
−w1 + w2 =0
w1 − w2 =0
which has the solution, w1 = 1 and w2 = 1. It is√common √ to normalize the eigenvector,
hence the eigenvector associated with λ = 1 is (1/ 2, 1/ 2). Repeating
√ the
√ procedure, we
should verify that the eigenvector associated with λ = −1 is (−1/ 2, 1/ 2). Finally, the
solution for the system linearized about the origin is:
√ √
1/√2 −1/√ 2
u
= c1 et + c2 e−t
v 1/ 2 1/ 2
where c1 and c2 are constants that cannot be found without initial conditions for our
linearized system.
Now let us examine the solution from Example 13.2. What does it say about
√ trajectories
√
that start near the origin? We see that along the direction of w1 = (1/ 2, 1/ 2), there
is an exponential
√ √growth away from the origin, λ > 0. However, along the direction of
w2 = (−1/ 2, 1/ 2), there is exponential decay towards the origin. This fixed point has
two different stabilities which depend on direction. The unstable direction w1 is tangent to
a curve called the unstable manifold of the fixed point. The stable direction w2 is tangent
to a curve called the stable manifold. Each vector w1 and w2 , is tangent to its respective
manifold at the fixed point. The origin in Example 13.2 is called a saddle point (or simply,
saddle) and is ultimately unstable, as we will see below. To better understand what is going
on here, let us examine the phase portrait near the origin.
Example 13.3: The phase portrait of the system from Example 13.2
Plot the phase portrait near the origin for the system of equations in Example 13.2.
Solution:
We used Mathematica to plot the phase portrait of the system. The code is shown in
Algorithm 13. Notice that there are four blocks of code. The first block simply defines
tmax the maximum time for which to integrate the equations. The second block finds
the solution to the system of equations for four different initial conditions, each of these
solutions is a trajectory in the phase portrait. The third block of code graphs the individual
parts of the phase portrait. The variables stable and unstable contain plots of the stable
and unstable manifold (locally near the fixed point), respectively. The formula for the
manifolds are obtained locally by using the point-point formula for a line using the points
w1 , w2 , and the origin. The variable stream contains a stream plot, which plots the local
direction of the vector field defined by the system of equations. The stream plot can be
a useful way of visualizing the flow of trajectories, as each arrow will be tangent to a
trajectory (similar to a velocity vector field!). The last four variables in the third block
of code store plots of each trajectory. Finally, the fourth block of code creates the graph
shown below.
i i
i i
i i
0.10
0.05
0.00
y
Out[ ]=
-0.05
-0.10
The black lines (in e-book) in the graph represent the manifolds and the red lines (in
e-book) are trajectories. The reader should try to identify each trajectory and manifold in
the graph with lines of code above. Notice that even though one direction of the saddle
is stable, i.e., attracting trajectories, the saddle is ultimately unstable because once the
trajectory gets close enough to the fixed point, the unstable direction begins to dominate
and trajectories run away from the saddle. Where do the trajectories go? We would need
to plot more of the phase portrait to find out. The reader should try that, too!
As we can see in the above example, the eigenvalues give the stability of the fixed
point. There are three simple types of fixed points associated with real eigenvalues in two-
dimensional systems: stable node, unstable node, and a saddle.
• The stable node has two negative eigenvalues and trajectories that start near the stable
node and approach it exponentially. A stable node corresponds to a stable equilibrium
position. For each stable node there is a region in the phase space called its basin of
attraction. The basin of attraction for a stable node is the set of all initial conditions
whose trajectories approach the stable node as t → ∞.
• The unstable node, sometimes called a repellor, has two positive eigenvalues. It is
called a repellor because trajectories starting near the node move away from the node
at an exponential rate. An unstable node corresponds to an unstable equilibrium.
• The saddle, explored above, has one positive and one negative eigenvalue. It corre-
sponds to an unstable equilibrium. The stable manifolds of the saddle often serve as
borders for the basin of attraction for stable nodes and other stable structures in the
phase space (see below).
In addition to the points listed above, there are three types of fixed points with complex
eigenvalues. Suppose λ = α + iβ, then eλt = eαt (cos(βt) + isin(βt)). In the case of complex
eigenvalues we see that there is a combination of exponential growth or decay and oscillation,
similar to the damped harmonic oscillator case. The possible fixed points are:
• A stable spiral occurs when Re(λ) < 0 and the trajectory spirals onto the fixed point.
A stable spiral appears in the phase portrait of the underdamped harmonic oscillator
with β < ω0 .
i i
i i
i i
tmax = 6;
soln1 = NDSolve[{x0 [t] == y[t] ∗ (x[t] + 1), y 0 [t] == x[t] ∗ (y[t] + 1), x[0] == −0.1, y[0] ==
0.099}, {x, y}, {t, 0, tmax}];
soln2 = NDSolve[{x0 [t] == y[t] ∗ (x[t] + 1), y 0 [t] == x[t] ∗ (y[t] + 1), x[0] == 0.1, y[0] ==
−0.101}, {x, y}, {t, 0, tmax}];
soln3 = NDSolve[{x0 [t] == y[t] ∗ (x[t] + 1), y 0 [t] == x[t] ∗ (y[t] + 1), x[0] == −0.1, y[0] ==
0.11}, {x, y}, {t, 0, tmax}];
soln4 = NDSolve[{x0 [t] == y[t] ∗ (x[t] + 1), y 0 [t] == x[t] ∗ (y[t] + 1), x[0] == 0.1, y[0] ==
−0.09}, {x, y}, {t, 0, tmax}];
i i
i i
i i
• An unstable spiral occurs when Re(λ) > 0 and the trajectory spirals away from the
fixed point. This would correspond to the damped harmonic oscillator with a negative
damping.
• A center occurs when Re(λ) = 0 and there is neither exponential growth away from the
fixed point, nor is there exponential decay onto the fixed point. The point is neutrally
stable and the trajectory circles the center. This corresponds to the undamped simple
harmonic oscillator.
What happens when there is a zero eigenvalue? One typically uses numerical methods to
analyze the system. Finally, it should be mentioned that there are cases that are called
degenerate because they have only one eigenvalue and/or one eigenvector. These lead to
other types of fixed points which are beyond the scope of this book. The interested reader
is directed to [Strogatz(2014)] for a thorough and readable review of even more types of
fixed points. Next, we will use our newfound knowledge of fixed points to analyze the simple
plane pendulum.
θ ℓ
The equation of motion of the pendulum has been found many times in this book,
θ̇ =y
(13.3.13)
ẏ = − ω02 sin θ
where the variable y is the angular velocity. Notice that fixed points occur at (θ∗ = 0, y = 0)
and (θ∗ = ±π, y = 0), these make sense physically because they correspond to the very
bottom and top of the pendulum’s path. Based on the way we have defined our coordinate
system, θ = ±π are the same point, the very top of the pendulum’s path. Note that θ = +π
when the pendulum approaches the top moving counterclockwise, and θ = −π when it
approaches the top moving clockwise. The Jacobian of the system is:
i i
i i
i i
y 0
-1
-2
-3
-3 -2 -1 0 1 2 3
θ
Figure 13.3: The phase portrait for (13.3.12) with ω0 = 1.
0 1
A= (13.3.14)
−ω02 cos θ 0
At the origin, the Jacobian (13.3.14) becomes:
0 1
(13.3.15)
−ω02 0
with eigenvalues λ = ±iω0 . The origin, corresponding to the bottom of the pendulum’s
path, is a center. The center is representing the oscillation of the pendulum about the
stable equilibrium position.
At (±π, 0),the Jacobian (13.3.14) becomes
0 1
(13.3.16)
ω02 0
which has eigenvalues λ = ±ω0 . The fixed point corresponding to the unstable equilibrium
is a saddle. Were you expecting an unstable node? The topmost point cannot be an unstable
mode because an unstable mode repels all trajectories. A saddle, on the other hand, has a
stable direction. We know that the pendulum can approach the top before turning around,
and hence, the topmost point must be a saddle. Now that we know the location and stability
of the fixed points, we are ready to draw the phase portrait.
Figure 13.3 illustrates the phase portrait for the simple pendulum using ω0 = 1. Recall
that the points θ = ±π correspond to the same physical location. Hence, the two saddle
points in the phase portrait are actually the same point. The red line in Figure 13.3 is both
the stable and unstable manifold of the saddle, which is sometimes referred to as a homoclinic
orbit, a trajectory that connects a saddle point to itself. Physically, the homoclinic orbit
represents the pendulum starting with an angular displacement of θ = π (vertical position)
and an angular velocity of 0, where the pendulum then swings down in a clockwise direction
through the stable equilibrium and returns to the vertical position, coming to a momentary
i i
i i
i i
1
θ (rad)
0
Out[ ]=
-1
θ(0) = 0.5
-2 θ(0) = 1.5
θ(0) = 2.5
0 2 4 6 8 10 12
t (sec)
Figure 13.4: The graphs of θ (in radians) versus time (seconds) for each of the oscillatory
trajectories shown in Figure 13.3.
rest and swinging back in a clockwise direction. The homoclinic orbit represents the largest
possible amplitude of oscillation of the pendulum.
Initial conditions that start inside the homoclinic orbit correspond to oscillatory behav-
ior. Notice that the trajectories (black lines) in Figure 13.3 form closed loops, which we know
to mean that the system eventually returns to its initial state after a time T , the period of
oscillation. Notice how the trajectories begin to take the shape of the homoclinic orbit as the
amplitude of oscillation gets bigger. The graphs of θ(t) for each of the trajectories inside the
homoclinic orbit are shown in Figure 13.4. Notice that as the amplitude increases in size,
the period of oscillation increases. Recall that in the simple harmonic oscillator, the period
of oscillation was independent of the amplitude of oscillation. In nonlinear oscillators, there
is typically a relationship between amplitude and period.
Finally, notice that the trajectories outside of the homoclinic orbit in Figure 13.3 do
not form closed loops. Those orbits correspond to the pendulum going over the top of its
motion and continually rotating in one direction, either clockwise (trajectories below the
homoclinic orbit) or counterclockwise (trajectories above the homoclinic orbit). Hence, we
see that the homoclinic orbit separates two distinct behaviors of the pendulum (oscillatory
versus continual rotation in one direction). The homoclinic orbit is sometimes called a
separatrix because it separates two distinct types of behavior.
i i
i i
i i
0.1
0.0
V(x)
Out[ ]=
-0.1
-0.2
ẋ =y
)
(13.3.19)
ẏ =x − x3
We see from (13.3.19) that fixed points exist at (0, 0), (1, 0), and (−1, 0). The Jacobian
matrix is,
0 1
A= (13.3.20)
1 − 3x2 0
√
which provides eigenvalues of ±1 for (0, 0) and ±i 2 for the points (1, 0) and (−1, 0).
Hence, the origin is a saddle and the other two points are centers. The phase portrait
can be plotted using a technique similar to that done in Example 13.3 and is shown in
Figure 13.6. Notice the two homoclinic orbits, one in red (in e-book) and the other in green
(in e-book), which surround the two centers. Inside each homoclinic orbit are closed-loop
trajectories corresponding to oscillations about each stable equilibrium point. Outside of
the homoclinic orbit are larger closed-loop trajectories which correspond to an oscillation
that passes through each equilibrium point.
The Mathematica code used to create Figure 13.6 is shown in Algorithm 14. The variable
names identify which part of the plot the particular line creates. The overlay of the stream
plot and the phase portrait trajectories can be very helpful when visualizing regions of the
phase portrait where we have not created trajectories. However, as the next subsection will
show, the overlay can make the phase portrait too “busy” and therefore difficult to read.
i i
i i
i i
0
y
Out[ ]=
-1
-2
-2 -1 0 1 2
x
Figure 13.6: The phase portrait for the double-well potential.
i i
i i
i i
ẋ =y
)
(13.3.22)
ẏ =x − x3 − 0.1y
We can see that (13.3.22) has three fixed points, (0, 0) and (±1, 0). The Jacobian matrix
of (13.3.22) is,
0 1
A= (13.3.23)
1 − 3x2 −0.1
and has eigenvalues of ±1 at the origin and −0.05 ± 1.41i for the other two points. Hence,
we see that the origin is a saddle and the other two points are a stable spiral. This should
not be surprising: the damping is causing decaying oscillation onto the equilibria at (±1, 0).
To better see this, we will study the phase portrait.
We create the phase portrait in Python using Algorithm 15. The equations of motion
are numerically solved using Python’s odeint command as we have done previously in this
text. There are four solutions created, one for each stable manifold and one for each unsta-
ble manifold. Notice that the stable manifolds, with variable names stable 1 and stable 2,
are found by integrating (13.3.22) backwards in time. The stable manifold is the set of
points that moves towards the saddle as t → ∞. The stable manifold can be difficult to
find integrating forward in time. Hence, to find the stable manifold, one chooses an initial
condition near the saddle point and integrates the equations backwards in time. Points on
the stable manifold moves away from the saddle as t → −∞. The unstable manifold is the
set of points that moves towards the saddle as t → −∞ (or away as t → ∞) is found by
solving the equations forward in time, using an initial condition near the saddle.
The phase portrait that results from Algorithm 15 is shown in Figure 13.7. In Figure
13.7, the saddle is represented by the circle at the origin. The unstable manifolds of the
saddle are shown in blue (in e-book) and each one connects to a stable spiral represented
by the squares. In red and green (in e-book) are the stable manifolds of the saddle. Notice
that they spiral around the saddle. Careful study of Figure 13.7 will show that the stable
manifolds define the basin of attraction for each stable spiral. Place the tip of a pencil
anywhere on the plot in between the red and green curves. From there follow the flow of the
solution by moving the pencil tip around the plot without crossing either the red or green
curve. Notice that the pencil tip will end up in one of the spirals. Next, repeat the process
but start in a different region, notice that the pencil may end up in another stable spiral!
The basins of attraction for these spirals wrap around each other. The initial condition you
choose will determine which spiral the pencil ends on. In other words, the final equilibrium
of the system is determined solely by the initial condition!
It is possible to create a stream plot in Python using the command streamplot found in
the matplotlib library. The code for a stream plot is shown in Algorithm 16. The code uses
the command mgrid from the Numpy library to set up a grid to compute the vector field.
The j in the mgrid command defines the number of steps in the grid. Hence, the grid runs
from -2 to 2 in each direction and that range is broken up into 1000 steps. The resulting
graph is shown in Figure 13.8, in which the direction of the flow is clearly seen. We chose not
to overlay the trajectories on the stream plot because the result was difficult to interpret.
i i
i i
i i
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
from s c i p y . i n t e g r a t e import o d e i n t
d e f f (X, t ) : #d e f i n e v e c t o r f i e l d
x1 , x2 = X
r e t u r n [ x2 , x1 − x1 ∗∗3 −0.1∗ x2 ]
time = np . l i n s p a c e ( 0 , 1 0 0 , 2 0 0 0 )
x0 = [ 0 . 0 1 , 0 . 0 ]
u n s t a b l e 1 = o d e i n t ( f , x0 , time )
p l t . p l o t ( u n s t a b l e 1 [ : , 0 ] , u n s t a b l e 1 [ : , 1 ] , ’ b− ’)
plt . plot ( [ unstable 1 [ 0 , 0 ] ] , [ unstable 1 [ 0 , 1 ] ] , ’o ’ )
plt . plot ( [ unstable 1 [ −1 ,0]] , [ unstable 1 [ −1 ,1]] , ’ s ’ )
x0 = [ − 0 . 0 1 , 0 . 0 ]
u n s t a b l e 2 = o d e i n t ( f , x0 , time )
p l t . p l o t ( u n s t a b l e 2 [ : , 0 ] , u n s t a b l e 2 [ : , 1 ] , ’ b− ’)
plt . plot ( [ unstable 2 [ −1 ,0]] , [ unstable 2 [ −1 ,1]] , ’ s ’ )
b a c k t i m e = np . l i n s p a c e ( 0 , −20 ,1 000)
x0 = [ 0 . 0 1 , 0 . 0 ]
s t a b l e 1 = o d e i n t ( f , x0 , b a c k t i m e )
plt . plot ( stable 1 [ : , 0 ] , stable 1 [ : , 1 ] , ’ g − ’)
x0 = [ − 0 . 0 1 , 0 . 0 ]
s t a b l e 2 = o d e i n t ( f , x0 , b a c k t i m e )
plt . plot ( stable 2 [ : , 0 ] , stable 2 [ : , 1 ] , ’ r − ’)
Algorithm 15: Python code to integrate (13.3.22) and plot the phase portrait.
import numpy a s np
import m a t p l o t l i b . p y p l o t a s p l t
w = 2
y , x = np . mgrid [−w:w: 1 0 0 0 j , −w : w: 1 0 0 0 j ]
x dot = y
y d o t = x− x ∗∗3 − 0 . 1 ∗ y
p l t . s t r e a m p l o t ( x , y , x dot , y d o t )
plt . xlabel ( ’x ’)
plt . ylabel ( ’y ’)
i i
i i
i i
i i
i i
i i
ẋ =y
)
(13.3.24)
ẏ =kx − x3 − by
The fixed points in (13.3.24) occur at the origin and at,
r !
k
(x∗ , y ∗ ) = ± ,0 (13.3.25)
If we think of k and as parameters that can be varied, then we see that the location
of the fixed points in (13.3.25) can change. For example, as k → 0 the fixed points (13.3.25)
collide at the origin. This collision is an example of a bifurcation and it tells us the condition
of which two stable equilibria can exist in the system. In this case, we use the term stable
equilibrium in the sense of Chapter 5, a local minimum in the potential. Furthermore, the
eigenvalues of the Jacobian matrix evaluated at (13.3.25) are,
1 p
λ= −b ± b2 − 8k (13.3.26)
2
When b > 0, we see that the fixed points located at (13.3.25) are stable spirals. However,
if b → 0, the previously fixed points become centers. This is an example of another type of
bifurcation where the fixed point changes stability. There are many types of bifurcations
between fixed points and a complete overview of them is beyond the scope of this book.
However, the interested reader should consult [Strogatz(2014), Hilborn(2001), Enns(2011)].
In the next paragraph, we will briefly mention two important bifurcations.
One of the most common bifurcations is the so-called saddle-node bifurcation. In a
saddle-node bifurcation, a saddle point and a stable node collide and both points disappear
from the phase portrait. Physically speaking, the system goes from having two equilibria,
one stable and one unstable, to having no equilibrium points (sometimes there is another
attracting structure elsewhere in the phase portrait). Another common bifurcation of fixed
points is the so-called transcritical bifurcation. In a transcritical bifurcation, a saddle and a
stable node collide, and pass through each other. However, as the collision occurs, the fixed
points exchange stabilities, so the saddle becomes a stable node and the stable node becomes
a saddle. Physically speaking, the equilibria exchange stabilities, but both equilibrium states
continue to exist after the bifurcation. Of course, for each bifurcation, the equilibrium state
is changing because the fixed point is moving, therefore changing the position and speed of
the equilibrium state.
We might imagine, that if we were designing a system we would want to know of possible
bifurcations in your system in order to prevent problems such as a desirable stable equi-
librium suddenly becoming unstable. For example, we might want to engineer an airplane
wing that stays relatively flat during flight. All airplane wings move a little during flight.
However, we would want to avoid the flat state from becoming unstable and the system
undergoing galloping, where the amplitude of oscillation of the wing gets large. A bifurca-
tion in which the galloping state becomes stable would be a dangerous situation, as can be
seen in videos of the famous Tacoma Narrows Bridge collapse.
i i
i i
i i
Next, we are going to study more complicated behaviors unique to nonlinear systems,
namely limit cycles and chaos.
1 d2 V 1 d3 V 1 d4 V
dV
V (x) ≈ V (0) + x+ x +
2
x +
3
x4 + · · · (13.4.1)
dx x=0 2! dx2 x=0 3! dx3 x=0 4! dx4 x=0
We have already made the argument that dV dx = 0 at the equilibrium, and that argument
continues to hold for the large amplitude case. Next, we know that the derivatives produce
a constant value when evaluated at the origin, so we will rewrite (13.4.1) in the form:
1 2 1 3 1 4
V (x) ≈ V (0) + kx + αx + δx + · · · (13.4.2)
2! 3! 4!
Next, we apply a restriction to the force which is common in physics, even in nonlinear
systems, that V (x) is symmetric about the equilibrium. Hence, the potential energy of the
object at a position x is the same as that at a position −x. In other words, we would not
expect the oscillator to have an energy that is dependent upon which side of equilibrium it
is on. In order to satisfy this restriction, we must require that α = 0. We can then find F
by calculating F = −dV /dx to get the force:
δ
F (x) = −kx − x3 (13.4.3)
6
We can then define a new parameter, = δ/6. The force becomes:
i i
i i
i i
0.8 1.4
d = 0.68
0.6 1.2
0.4 1.0
0.2 0.8
x'
x
Out[ ]=
0.0 0.6
-0.2 0.4
-0.4 0.2
-0.6 0.0
0.0 0.2 0.4 0.6 0.8 1.0 1.2 1.4 4500 4520 4540 4560 4580 4600
x t
Figure 13.9: The phase portrait (left) and x(t) (right) of (13.4.6) with ci = 1, for all i, ω = 1,
and d = 0.68.
ẍ + c1 ẋ − c2 x + c3 x3 = d cos(ωt) (13.4.6)
where c1 = β/m, c2 = k/m, c3 = /m, and d = f /m. Equation (13.4.6) is called the Duffing
equation and is a standard model equation used in nonlinear dynamics to study a wide range
of nonlinear phenomena. With the right choice of parameters, ci and d, the Duffing equation
can be made to display a large variety of nonlinear behaviors. Because the behavior of the
Duffing equation is so well known, it is often used as a model system to test new algorithms
for analyzing nonlinear systems. In the next subsections, we will use the Duffing equation
to demonstrate limit cycles, period-doubling bifurcations, and chaos.
i i
i i
i i
0.8 1.4
d = 0.69
0.6 1.2
0.4 1.0
0.2 0.8
x'
x
Out[ ]=
0.0 0.6
-0.2 0.4
-0.4 0.2
-0.6 0.0
0.0 0.2 0.4 0.6 0.8 1.0 1.2 1.4 4500 4520 4540 4560 4580 4600
x t
Figure 13.10: The phase portrait (left) and x(t) (right) of (13.4.6) with ci = 1, for all i,
ω = 1, and d = 0.69.
1.5
0.8
d = 0.75
0.6
1.0
0.4
0.2
x'
Figure 13.11: The phase portrait (left) and x(t) (right) of (13.4.6) with ci = 1, for all i,
ω = 1, and d = 0.75.
dimensions. Although beyond the scope of this Chapter, the Duffing equation is actually a
three-dimensional system due to the explicit dependence of time in this system. The solution
for d = 0.69 is called a period-2 solution because, loosely speaking, the oscillator has to go
through two oscillations before returning to its initial state, once the oscillator has settled
onto the limit cycle behavior.
Figure 13.11 shows the phase portrait (left) and solution (right) when d is increased
to 0.75. There are even more apparent intersections in the phase portrait and the solution
x(t) appears even more complicated. Figure 13.11 illustrated a period-4 solution to (13.4.6).
We are seeing an interesting phenomenon that can occur in nonlinear oscillators called a
period-doubling bifurcation, whereas the parameter d in (13.4.6) varies, the period of the
oscillator (and therefore the limit cycle) doubles. One method of illustrating period-doubling
bifurcations is by using something called a Poincaré section, which we will discuss below.
The Poincaré section samples the solution once a period and then plots the phase space.
Since the drive frequency is ω = 1, we can sample x(t) and ẋ(t) at times t = 2nπ, where
n is an integer. The code for producing Poincaré sections is shown below in Algorithm 17.
The variable, psdata, stores a table of sampled solutions, every n2π timesteps. The First
command is used because Mathematica introduces some extra curly brackets when this
method of solution sampling is used.
i i
i i
i i
0.75 0.75
d = 0.68 d = 0.69
0.5 0.5
0.25 0.25
x'
x'
0 0
-0.25 -0.25
-0.5 ● -0.5 ●
●
0 0.3 0.6 0.9 1.2 0 0.3 0.6 0.9 1.2
x x
Out[ ]=
0.75
d = 0.75
0.5
0.25
x'
●
0
-0.25 ●
●
-0.5 ●
0 0.3 0.6 0.9 1.2
x
Figure 13.12: The Poincaré sections for the solutions shown in Figures 13.9 through 13.11.
The Poincaré sections for the solutions in Figures 13.9 through 13.11 are shown in Figure
13.12.
Notice how in Figure 13.12, for d = 0.68, there is only one point. That is because, after
settling onto the limit cycle, the solution repeats its state once every period. Notice that we
are not sampling the solution at its maximum; the x-coordinate of the point in the Poincaré
section does not correspond to the maximum value of the top right graph in Figure 13.9.
The fact that the point in the Poincaré section does not match either the maximum or
minimum value of the solution does not matter. What matters is that every t = 2nπ, the
system returns to the same state. However, we see that when d = 0.69, the oscillator appears
to jump between two points in its Poincaré section. It takes 4π = 2 × 2π, or twice the period
to repeat its state. Continuing this reasoning, we see that when d = 0.75, the system is
exhibiting period-4 behavior.
As mentioned previously, Figures 13.9 through 13.12 are demonstrating period-doubling
bifurcations. We see a period-doubling bifurcation in the limit cycle between d = 0.68 and
d = 0.69. We see another period-doubling bifurcation between d = 0.69 and d = 0.75. As we
continue to increase d, we would see several more period doubling bifurcations, the series
of which is called a period-doubling cascade, which will eventually lead to chaotic behavior
i i
i i
i i
Phase Space
x(t)
1.0
1.5
0.5 1.0
0.5
x'
0.0 0.0
x
Out[ ]=
-0.5
-0.5 -1.0
-1.5
-1.0
4500 4600 4700 4800 4900 5000
-1.5 -1.0 -0.5 0.0 0.5 1.0 1.5 t
x
Poincaré Section
0.5
0.0
x'
Out[ ]=
-0.5
-1.0
-1.0 -0.5 0.0 0.5 1.0 1.5
x
Figure 13.13: Plots of the phase portrait, solution, and Poincaré section of (13.5.1).
once d = 0.8. However, to illustrate chaos, we will change the parameters used in (13.4.6)
in order to better visualize the attractor corresponding to chaotic behavior.
13.5 CHAOS
In addition to interesting oscillations, nonlinear systems can also display a behavior called
chaos. To explore the phenomenon of chaos, we will use the Duffing equation, but with a
specific set of parameter values:
i i
i i
i i
that the Poincaré section is not a simple collection of points. In fact, it has its own complex
structure. You might wonder if we solved (13.5.1) for a longer period of time, would we find
that any of the points in the Poincaré section repeat themselves, hence demonstrating a
periodicity. The answer is no, there is no repetition, and therefore the system is not periodic.
The Poincaré section is even more interesting if we zoom into it and look at its structure.
Doing so would reveal that the bands that are visible in the Poincaré section of Figure
(13.13) persist when zoomed in at any scale provided enough points are plotted, which can
get computationally expensive. The presence of structure at arbitrary scale means that the
attractor is a fractal!
So what have we learned in our study so far from Figure (13.13)? We learned that chaotic
systems are not periodic. The trajectory of a chaotic system’s solution in the phase portrait
attracts to a structure with fractal properties. Furthermore, the behavior of the solution is
not random; the system is obeying the Duffing equation. This is getting strange! But wait,
there is more...
i i
i i
i i
x
-0.5
-0.5
-1.0
-1.0 -1.5
0 10 20 30 40 50 0 20 40 60 80 100
t t
Out[ ]=
x
-0.5 -0.5
-1.0 -1.0
-1.5 -1.5
0 20 40 60 80 100 0 20 40 60 80 100
t t
Figure 13.14: Plots comparing similar initial conditions for the simple harmonic oscillator
(SHO) and the Duffing equation. The value of Difference in the title of each plot descibes
the difference in the inital positions of the two oscillators. The red (in e-book) curve has
an initial position of x(0) = 0 and the blue (in e-book) curve has an initial position of
x(0) = 0 + Difference. In all cases ẋ(0) = 0.
Attractors on which the trajectories exhibit sensitive dependence on initial conditions are
called strange attractors.
Is the sensitive dependence on initial conditions just a fancy way of saying that the
system is random? No, the system is not random. The Duffing equation is an ordinary
differential equation. The solution of the Duffing equation is dictated by its initial conditions,
and is therefore, a deterministic system. In other words, future values of x(t) and ẋ(t) depend
on the initial conditions x(0) and ẋ(0). In a random system, past states do not dictate future
states (otherwise there would be no lottery!). The fact that this system is deterministic
means that the interesting behavior arises due to the nonlinearity of the system, not some
random elements.
Other than being a mathematical curiosity, why is sensitive dependence on initial con-
ditions important? As an example, it is helpful to consider the case where chaos was first
presented in its modern form by Edward Lorenz, a meteorologist, who was creating weather
models. Lorenz’s paper [Lorenz(1963)] is very readable, consider checking it out. Lorenz’s
weather model demonstrated chaotic behavior, so as an example, let us suppose that the
weather perfectly follows Lorenz’s model (in reality, it does not). Now, we are excited
because we have the perfect weather model and therefore should be able to predict the
weather at any time in the future. We are ready for a lucrative future as a weather fore-
caster, but not so fast—we realize that in order to make predictions, we need to input
an initial condition into the model. The initial condition comes from instruments used to
measure various atmospheric conditions. The instruments have a certain degree of accuracy,
they may only measure the true value of a quantity to for example, one part in a thousand.
The actual degree of accuracy is irrelevant, the point is that the measured initial condi-
tions will be different from the true state of the atmosphere. Since, in our example, the
weather perfectly obeys a chaotic model, then the slight difference in measured versus true
i i
i i
i i
values will mean that after a period of time, the model’s predictions will diverge from the
actual weather. The chaotic system is sensitive to initial conditions! The result of a chaotic
system’s sensitive dependence on initial conditions is that there is no accurate long-term
prediction possible for chaotic systems.
i i
i i
i i
In this Chapter, we discussed that nonlinear systems can take on a wide variety of
behaviors. To be honest, we have not even scratched the surface on the variety of interest-
ing phenomena that can occur in nonlinear systems. As promised, we will include a few extra
suggestions on how to extend the reader’s study of nonlinear systems. We recommend start-
ing with [Strogatz(2014)], as the book is easily understandable and enjoyable to read. From
there, one can take many pathways to learning nonlinear systems. A book we frequently
recommend to students is [Enns(2011)], which contains some very interesting applications of
nonlinear systems to a wide variety of fields in the natural and social sciences. Another book
commonly used in nonlinear dynamics courses is [Hilborn(2001)]. If the reader is interested
in analyzing data measured from nonlinear systems, both [Kantz and Schreiber(2004)] and
[Abarbanel(1996)] are recommended. Finally, if the rader is interested in developing models
of systems (especially complex systems), we recommend [Wilensky and Rand(2015)].
If one spends enough time working with (and learning about) nonlinear systems, even-
tually, one will come to start thinking of problems not as linear versus nonlinear, or even as
physics versus not physics. Instead, one will likely just see the system you are working with
as just that, an interesting system to be analyzed with one of many tools the reader will
have in his/her mathematical and computational toolbox. In the end, breaking up problems
in the natural world by discipline is a human-imposed construct. Real systems are neither
biology versus physics versus sociology, they are just systems to be studied. While the con-
struct of silos for problems is useful for the purposes of education, once one no longers need
it, we recommend discarding it and focus on studying problems that are interesting to you.
ẋ =f (x, y)
ẏ =g(x, y)
Although the equations governing nonlinear systems are not usually solvable in closed
form, we can use tools such as phase portraits to understand possible steady state behaviors
of the system as well as identify the stability of equilibria.
The fixed points of the system can be found by solving
f (x∗ , y ∗ ) =0
g(x∗ , y ∗ ) =0
Although we cannot solve the whole system analytically, we can linearize the system
about the fixed points, which results in an easily solvable system of equations:
∂f ∂f
u̇ ∂x ∂y u
=
v̇ ∂g ∂g v
∂x ∂y
where the eigenvalues of the Jacobian matrix A give us information about the stability of
the equilibrium. The eigenvalues can be found using the characteristic polynomial,
det (A − λ1) = 0
i i
i i
i i
In nonlinear systems, the equilibrium states can change (both location in phase space and
stability).
Finally, we also explored a behavior called chaos, which is unique to nonlinear systems.
A hallmark of chaos is that the solutions of a chaotic system display sensitive dependence
on initial conditions. In other words, two initially similar states will evolve to very different
behaviors, thus inhibiting one’s ability to make long-term predictions on chaotic systems.
d3 x
+x = 0
dt3
Is this equation a linear or a nonlinear differential equation? Justify your answer. Solve
the equation for x(t).
2. Consider the differential equation
d3 x
+ t3 x = 0
dt3
ẍ = t2 x − tx2
Suppose you have found two solutions to this equation, x1 (t) and x2 (t). Is it true that
c1 x1 + c2 x2 is also a solution? Prove your answer.
Section 13.2: The Damped Harmonic Oscillator, Revisited
4. Using a computer, plot the phase portrait for the damped harmonic oscillator with
γ < 0 and ω02 = 1. You may use x(0) = 1 and ẋ(0) = 0. Do you get what you expect?
Why or why not? What does it mean for γ < 0?
5. Using a computer, plot the phase portrait for the damped pendulum, ẍ + ω02 sin x +
γ ẋ = 0 for ω02 = 1 and three values of γ, less than, equal to, and greater than ω02 . For
each case use x(0) = 1 and ẋ(0) = 0. Discuss your results.
6. Repeat Problem 5 using x(0) = 3 and ẋ(0) = 0. Discuss your results and compare them
to that of Problem 5.
Section 13.3: Fixed Points and Phase Portraits
i i
i i
i i
which makes the potential function similar to the potential energy discussed in Chap-
ter 5. In this chapter, we discussed that the condition for fixed points (i.e., equilibrium
points), x∗ , f (x∗ ) = 0. How can we use the potential to find fixed points using V (x)?
In other words, how does one find fixed points using V (x)? How can one use V (x) to
identify the stability of a fixed point?
8. Consider the system ẋ = 1 − x2 . Find the potential corresponding to this system (see
Problem 7). Using the potential, find the fixed points and identify their stability.
9. In this chapter, we discussed how to identify the stability of fixed points for two-
dimensional systems. The logistic equation, ẋ = rx(1 − x/k), is used to model popula-
tion sizes, and is a one-dimensional system. Note that r is a positive constant called
the linear growth rate and k is a positive constant called the carrying capacity. Find
the fixed points for the logistic equation and give them a physical interpretation. Find
the stability of each fixed point by perturbing the fixed point, x = x∗ + u (for u << 1)
and expanding the logistic equation in powers of u. The result will be a differential
equation for u which you can solve.
10. Find the eigenvalues and eigenvectors for the following matrices:
0 1 2 1 1 1
(a) (b) (c)
3 0 7 3 1 1
Ṙ =aR + bJ
J˙ =cJ + dR
and all kinds of interesting relationships can be created by choosing different values
of the parameters a, b, c, and d. Following the methods of Section 13.3 solve the love
affair model in closed form for an eager Romeo (a = 1, b = 2) and a very cautious Juliet
(c = −1, d = −3). Plot the phase portrait and comment on the possible outcomes for
various initial conditions of the relationship.
14. In this problem, we will build on Problem 13 with a nonlinear romance! We will have
both Romeo and Juliet be very sensitive to their own emotions:
Ṙ = − 2R − 2J(1 − |J|)
J˙ =J + R(1 − |R|)
i i
i i
i i
The nonlinear term J(1 − |J|) is sometimes called a repair nonlinearity [Sprott(2004)]
and the number 1 in the nonlinear term is a measure of when Juliet’s love becomes
counter productive. Plot the phase portrait for this system and discuss the nature of
the relationship between Romeo and Juliet.
15. Find the Jacobian matrix for (13.3.24) and compute its eigenvalues for each of the
fixed points.
16. The Lotka-Volterra model of competition is a mathematical model for two species X
and Y who compete for the same resources. Let x be the population size of species X
and y be the population size of species Y , then one choice of a Lotka-Volterra model
parameters would be:
ẋ =4x − 3x2 − xy
ẏ =2y − 4y 2 − 3xy
For this model, identify the fixed points and their stability. Can the two species ever
coexist? Note that this kind of equation-based model tends to work well for large
populations. For small populations, agent-based models tend to be more useful. For a
discussion of agent-based versus equation-based modeling, the interested reader should
consult [Wilensky and Rand(2015)]
17. Plot the phase portrait for the Lotka-Volterra competition model in Problem 16.
Discuss the final state of the system for various initial conditions.
18. Lotka-Volterra equations can also be used to describe predator-prey models where
species Y (predator) eats species X (prey). Consider the predator-prey model
ẋ =3x − 7xy
ẏ = − 2y + 5xy
where like the competition model, we will assume x(t) and y(t) measure the number
of each species in the hundreds. Find the fixed points of this predator-prey model and
identify their stability. Comment on the nature of the equilibrium state corresponding
to each fixed point.
19. Draw the phase portrait for the predator-prey model in Problem 18.
20. The SIS (susceptible-infected-susceptible) disease transmission model can be used for
diseases that do not remove individuals from the population, either through immunity
or death. SIS models are sometimes appropriate for some bacterial diseases such as
streptococcal pharyngitis (also known as sore throat). Consider two populations of
individuals s(t) and i(t), representing the number of susceptible and infected individ-
uals respectively. The SIS model is:
ṡ =µn − µs + γi − βsi
i̇ = − µi − γi + βsi
where the coefficients represent rates and n = s + i is the total population. Comment
on the physical meaning of each coefficient: µ,γ, and β. Find the fixed points for the
system and identify their stability. Note that negative values of s and i are meaningless.
i i
i i
i i
Notice that the drive frequency (the time derivative of the drive phase) is fd = ω0 − αt,
so that at t = 0, the drive frequency is equal to the natural frequency of the pendulum.
Starting with initial conditions, x(−1000) = 0 and ẋ(−1000) = 0, compute x(t) for
= 0.0459 and = 0.0461 for the pendulum with ω0 = 2π, α = 0.001. For which value
of is the pendulum displaying autoresonance? For more information on autoresonance
and the pendulum, check out [Fajans and Friédland(2001)].
24. Electrical circuits provide excellent examples of nonlinear systems. One of the most
famous circuit-inspired equations is the so-called Van der Pol equation
ẍ − 1 − x2 ẋ + x = 0
which is a simple harmonic oscillator with a nonlinear damping term. Notice that the
sign of the damping actually changes with the value of x. For example if > 0, then
the damping is negative for x < 1 and is positive for x > 1. Thus small oscillations
will grow due to negative damping, but once the oscillations are large enough, they
will dampen out as the damping term becomes positive. The Van der Pol oscillator
can also display so-called relaxation oscillations, a common form of oscillations in
nonlinear systems where rapid changes of x are followed by slower variations. To see
relaxation oscillations in action, solve the Van der Pol oscillator for = 8.0. Plot both
x(t) and the phase portrait. Describe the solution.
25. It can sometimes be convenient to describe limit cycles using polar coordinates.
Consider the following system in polar coordinates:
ṙ =r(r − 1)(r − 2)
θ̇ =1
i i
i i
i i
where r and θ are the usual polar coordinates with r ≥ 0. Notice that one can find
values of r that make ṙ = 0, but there are no values of θ or r for which θ̇ = 0. That
is OK! Give a description of the phase portrait in this case. There are limit cycles in
this system. Where are they and what is their stability? The stability can be found
by using the method outlined in Problem 9.
Section 13.5: Chaos
26. The Lorenz equations:
ẋ =σ(y − x)
ẏ =rx − y − xz
ż =xy − bz
display chaotic behavior with σ = 10, r = 28, and b = 8/3. Plot the phase portrait (x,
y, z) for the Lorenz equations.
27. The Lorenz equations in Problem 26 are a chaotic system. Show that they demonstrate
sensitive dependence on initial conditions.
28. Chua’s circuit is a electrical circuit with a piecewise linear negative resistance NR . It
is one of the simplest circuits to exhibit a variety of nonlinear behaviors including a
period-doubling route to chaos. The circuit schematic is shown in Figure 13.15. Let
x, y, and z be proportional to the voltages across the capacitors C1 and C2 , and the
current in the inductor L, respectively. The equations governing the circuit can be
obtained from Kirchoff’s rules and simplified to be:
ẋ =α(y − x − g(x))
ẏ =x − y + z
ż = − βy
R
NR
L C1 C2
29. A period-doubling route to chaos can be observed in the equations for Chua’s circuit
(see Problem 28) by using the same parameters as used in Problem 28, but with
β = 50, 35, 33.8, etc. For each value of β, create a graph of x(t) and identify the
periodicity of each solution. For what value of β is there a period-8 solution?
i i
i i
i i
ẋ =x(1 − x) − xy
ẏ = − y + xy
i i
i i
i i
Bibliography
445
i i
i i
i i
446 Bibliography
[Kantz and Schreiber(2004)] H Kantz and T Schreiber. Nonlinear Time Series Analysis.
Cambridge University Press, Cambridge, UK, 2nd edition, 2004.
[Kinder and Nelson(2015)] JM Kinder and P Nelson. A Student’s Guide to Python for
Physical Modeling. Princeton University Press, Princeton, NJ, 2015.
[Lorenz(1963)] E Lorenz. Deterministic nonperiodic flow. Journal of Atmospheric Sciences,
20:130 – 141, 1963.
[Morin(2008)] D Morin. Introduction to Classical Mechanics with Problems and Solutions.
Cambridge, Cambridge, 2008.
[Noether(1918)] E Noether. Invariante variationsprobleme. Nachrichten von der
Gesellschaft der Wissenschaften zu GÃűttingen, Mathematisch-Physikalische Klasse,
1918:235–257, 1918.
[Piessens et al. (1983)] R Piessens, E. de Doncker-Kapenga, C.W. Überhuber, and D.K.
Kahaner. QUADPACK: A Subroutine Package for Automatic Integration. Springer,
New York, 1983.
[Press et al. (2007)] WH Press, SA Teukolsky, WT Vetterling, and BP Flannery. Numerical
Recipes, 3rd Edition. Cambridge University Press, Cambridge, UK, 2007.
[Sprott(2004)] JC Sprott. Dynamical models of love. Nonlinear Dynamics, Psychology, and
Life Sciences, 8(8):303 – 313, 2004.
[Strogatz(2014)] S Strogatz. Nonlinear Dynamics and Chaos: With Applications to Physics,
Biology, Chemistry, and Engineering. Westview Press, Cambridge, UK, 2 edition, 2014.
[Supiano(2018)] B. Supiano. So what are you going to do with that degree? physics majors
get that question, too. The Chronicle of Higher Education, 9 2018.
[Taylor(2005)] JR Taylor. Classical Mechanics. University Science Books, Sausalito, CA,
2005.
[Thornton and Marion(2004)] ST Thornton and JB Marion. Classical Dynamics of Parti-
cles and Systems. Thomson-Brooks/Cole, Belmont, CA, 5 edition, 2004.
[Tokpah(2008)] C.L. Tokpah. The Effects of Computer Algebra Systems on Students’
Achievement in Mathemaics. dissertation, Kent State University, 2008.
[Wilensky and Rand(2015)] U Wilensky and W Rand. An Introduction to Agent-Based
Modeling. MIT Press, Cambridge, MA, 2015.
i i
i i
i i
Index
447
i i
i i
i i
448 Index
i i
i i
i i
Index 449
i i
i i
i i
450 Index
i i
i i
i i
Index 451
i i
i i
i i
452 Index
i i
i i
i i
Index 453
i i
i i
i i
454 Index
i i
i i
i i
Index 455
i i
i i