Mathematical Methods Notes
Mathematical Methods Notes
MATHEMATICAL METHODS
Mihir Sen
Joseph M. Powers
updated
28 March 2011, 10:32am
2
1 Multi-variable calculus 13
1.1 Implicit functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 13
1.2 Functional dependence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 16
1.3 Coordinate transformations . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
1.3.1 Jacobians and metric tensors . . . . . . . . . . . . . . . . . . . . . . . 21
1.3.2 Covariance and contravariance . . . . . . . . . . . . . . . . . . . . . . 28
1.4 Maxima and minima . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 36
1.4.1 Derivatives of integral expressions . . . . . . . . . . . . . . . . . . . . 37
1.4.2 Calculus of variations . . . . . . . . . . . . . . . . . . . . . . . . . . . 38
1.5 Lagrange multipliers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
Problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 46
3
4 CONTENTS
10 Appendix 415
10.1 Trigonometric relations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 415
10.2 Routh-Hurwitz criterion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 416
10.3 Infinite series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 417
10.4 Asymptotic expansions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 418
10.5 Special functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 418
10.5.1 Gamma function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 418
10.5.2 Beta function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 419
10.5.3 Riemann zeta function . . . . . . . . . . . . . . . . . . . . . . . . . . 419
10.5.4 Error function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 419
10.5.5 Fresnel integrals . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 420
10.5.6 Sine- and cosine-integral functions . . . . . . . . . . . . . . . . . . . . 420
10.5.7 Elliptic integrals . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 421
10.5.8 Gauss’s hypergeometric function . . . . . . . . . . . . . . . . . . . . . 422
10.5.9 δ distribution and Heaviside function . . . . . . . . . . . . . . . . . . 422
Bibliography 429
These are lecture notes for AME 60611 Mathematical Methods I, the first of a pair of courses
on applied mathematics taught at the Department of Aerospace and Mechanical Engineering
of the University of Notre Dame. Until Fall 2005, this class was numbered as AME 561. Most
of the students in this course are beginning graduate students in engineering coming from a
wide variety of backgrounds. The objective of the course is to provide a survey of a variety
of topics in applied mathematics, including multidimensional calculus, ordinary differential
equations, perturbation methods, vectors and tensors, linear analysis, and linear algebra, and
dynamic systems. The companion course, AME 60612, covers complex variables, integral
transforms, and partial differential equations.
These notes emphasize method and technique over rigor and completeness; the student
should call on textbooks and other reference materials. It should also be remembered that
practice is essential to the learning process; the student would do well to apply the techniques
presented here by working as many problems as possible.
The notes, along with much information on the course itself, can be found on the world
wide web at http://www.nd.edu/∼powers/ame.60611. At this stage, anyone is free to
duplicate the notes on their own printers.
These notes have appeared in various forms for the past few years; minor changes and
additions have been made and will continue to be made. Thanks especially to Prof. Bill
Goodwine and his Fall 2006 class who identified several small errors. We would be happy to
hear from you about further errors or suggestions for improvement.
Mihir Sen
[email protected]
http://www.nd.edu/∼msen
Joseph M. Powers
[email protected]
http://www.nd.edu/∼powers
The content of this book is licensed under Creative Commons Attribution-Noncommercial-No Derivative Works 3.0.
11
12 CONTENTS
Multi-variable calculus
if ∂f /∂y 6= 0.
The derivative ∂y/∂xi can be determined from f = 0 without explicitly solving for y.
First, from the chain rule, we have
∂f ∂f ∂f ∂f ∂f
df = dx1 + dx2 + . . . + dxi + . . . + dxn + dy = 0. (1.2)
∂x1 ∂x2 ∂xi ∂xn ∂y
Differentiating with respect to xi while holding all the other xj , j 6= i constant, we get
∂f ∂f ∂y
+ = 0, (1.3)
∂xi ∂y ∂xi
so that
∂f
∂y
= − ∂x
∂f
i
, (1.4)
∂xi ∂y
f (x, y, u, v) = 0, (1.5)
g(x, y, u, v) = 0. (1.6)
13
14 CHAPTER 1. MULTI-VARIABLE CALCULUS
Under certain circumstances, we can unravel these equations (either algebraically or numer-
ically) to form u = u(x, y), v = v(x, y). The conditions for the existence of such a functional
dependency can be found by differentiation of the original equations, for example:
∂f ∂f ∂f ∂f
df = dx + dy + du + dv = 0. (1.7)
∂x ∂y ∂u ∂v
Holding y constant and dividing by dx we get
∂f ∂f ∂u ∂f ∂v
+ + = 0. (1.8)
∂x ∂u ∂x ∂v ∂x
In the same manner, we get
∂g ∂g ∂u ∂g ∂v
+ + = 0, (1.9)
∂x ∂u ∂x ∂v ∂x
∂f ∂f ∂u ∂f ∂v
+ + = 0, (1.10)
∂y ∂u ∂y ∂v ∂y
∂g ∂g ∂u ∂g ∂v
+ + = 0. (1.11)
∂y ∂u ∂y ∂v ∂y
Equations (1.8,1.9) can be solved for ∂u/∂x and ∂v/∂x, and Eqs. (1.10,1.11) can be solved
for ∂u/∂y and ∂v/∂y by using Cramer’s1 rule. To solve for ∂u/∂x and ∂v/∂x, we first write
Eqs. (1.8,1.9) in matrix form:
∂f ∂f ∂u
∂u ∂v ∂x
− ∂f
∂x
∂g ∂g ∂v = ∂g . (1.12)
∂u ∂v ∂x − ∂x
∂f ∂f ∂f
− ∂y ∂v
∂u − ∂f∂y
∂g ∂g ∂(f,g) ∂g ∂g ∂(f,g)
∂u − ∂v − ∂y
∂(y,v) ∂(u,y)
= ∂f∂y ∂v
∂f ≡ − ∂(f,g) , ∂u
= ∂f ∂f ≡ − ∂(f,g) . (1.14)
∂y ∂u ∂v ∂y ∂u ∂v
∂(u,v) ∂(u,v)
∂g ∂g ∂g ∂g
∂u ∂v ∂u ∂v
1
Gabriel Cramer, 1704-1752, well-traveled Swiss-born mathematician who did enunciate his well known
rule, but was not the first to do so.
If the Jacobian2 determinant, defined below, is non-zero, the derivatives exist, and we
indeed can form u(x, y) and v(x, y).
∂f ∂f
∂(f, g) ∂u
6 0.
= ∂g ∂v
∂g = (1.15)
∂(u, v) ∂u ∂v
This is the condition for the implicit to explicit function conversion. Similar conditions hold
for multiple implicit functions fi (x1 , . . . , xn , y1 , . . . , ym ) = 0, i = 1, . . . , m. The derivatives
∂fi /∂xj , i = 1, . . . , m, j = 1, . . . , n exist in some region if the determinant of the matrix
∂fi /∂yj 6= 0 (i, j = 1, . . . , m) in this region.
Example 1.1
If
x + y + u6 + u + v = 0, (1.16)
xy + uv = 1. (1.17)
Find ∂u/∂x.
Note that we have four unknowns in two equations. In principle we could solve for u(x, y) and
v(x, y) and then determine all partial derivatives, such as the one desired. In practice this is not always
possible; for example, there is no general solution to sixth order equations such as we have here.
The two equations are rewritten as
f (x, y, u, v) = x + y + u6 + u + v = 0, (1.18)
g(x, y, u, v) = xy + uv − 1 = 0. (1.19)
Using the formula developed above to solve for the desired derivative, we get
∂f ∂f
−
∂x ∂g
∂v
∂g
∂u −
= ∂f∂x ∂f∂v . (1.20)
∂x ∂u ∂v
∂g ∂g
∂u ∂v
Substituting, we get
−1 1
∂u −y u y−u
= 5 =
5 + 1) − v
. (1.21)
∂x 6u + 1 1
u(6u
v u
Note when
v = 6u6 + u, (1.22)
that the relevant Jacobian is zero; at such points we can determine neither ∂u/∂x nor ∂u/∂y; thus we
cannot form u(x, y).
At points where the relevant Jacobian ∂(f, g)/∂(u, v) 6= 0, (which includes nearly all of the (x, y)
plane) given a local value of (x, y), we can use algebra to find a corresponding u and v, which may be
multivalued, and use the formula developed to find the local value of the partial derivative.
2
Carl Gustav Jacob Jacobi, 1804-1851, German/Prussian mathematician who used these determinants,
which were first studied by Cauchy, in his work on partial differential equations.
∂f ∂u ∂f ∂v
+ = 0, (1.23)
∂u ∂x ∂v ∂x
∂f ∂u ∂f ∂v
+ = 0, (1.24)
∂u ∂y ∂v ∂y
∂u ∂v ∂f
∂x ∂x ∂u 0
∂u ∂v ∂f = . (1.25)
∂y ∂y ∂v
0
Since the right hand side is zero, and we desire a non-trivial solution, the determinant of the
coefficient matrix, must be zero for functional dependency, i.e.
∂u ∂v
∂x
∂u ∂x = 0. (1.26)
∂v
∂y ∂y
Example 1.2
Determine if
u = y + z, (1.28)
v = x + 2z 2 , (1.29)
w = x − 4yz − 2y 2 , (1.30)
0 1 1
= 1 0 −4(y + z) ,
(1.32)
1 4z −4y
= (−1)(−4y − (−4)(y + z)) + (1)(4z), (1.33)
= 4y − 4y − 4z + 4z, (1.34)
= 0. (1.35)
Example 1.3
Let
x+y+z = 0, (1.36)
x + y + z 2 + 2xz
2 2
= 1. (1.37)
f (x, y, z) = x + y + z = 0, (1.38)
g(x, y, z) = x2 + y 2 + z 2 + 2xz − 1 = 0 (1.39)
∂f ∂f ∂f
df = dz + dx + dy = 0, (1.40)
∂z ∂x ∂y
∂g ∂g ∂g
dg = dz + dx + dy = 0, (1.41)
∂z ∂x ∂y
∂f ∂f dx ∂f dy
+ + = 0, (1.42)
∂z ∂x dz ∂y dz
∂g ∂g dx ∂g dy
+ + = 0, (1.43)
∂z ∂x dz ∂y dz
∂f ∂f dx ∂f
∂x ∂y dz − ∂z
∂g ∂g dy = ∂g , (1.44)
∂x ∂y dz − ∂z
T
then the solution matrix (dx/dz, dy/dz) can be obtained by Cramer’s rule:
∂f ∂f
− ∂z ∂y −1 1
∂g ∂g
dx − −(2z + 2x) 2y −2y + 2z + 2x
= ∂f∂z ∂f∂y = = = −1, (1.45)
dz ∂x ∂y 1 1
2y − 2x − 2z
∂g ∂g 2x + 2z 2y
∂x
∂f ∂y ∂f
∂x − ∂z
1 −1
∂g ∂g
dy − 2x + 2z −(2z + 2x) 0
= ∂x ∂f
∂z
∂f = = . (1.46)
dz ∂x ∂y
1 1
2y − 2x − 2z
∂g ∂g 2x + 2z 2y
∂x ∂y
Note here that in the expression for dx/dz that the numerator and denominator cancel; there is no
special condition defined by the Jacobian determinant of the denominator being zero. In the second,
dy/dz = 0 if y − x − z 6= 0, in which case this formula cannot give us the derivative.
2 -1 x
0
1
1
y
0
1
-1
0.5
-2 0 z
1
-0.5
0.5
z -1
0 0.5
-0.5 0 y
-1 -0.5
-1
-0.5
0
x 0.5
1
Now, in fact, it is easily shown by algebraic manipulations (which for more general functions are
not possible) that
√
2
x(z) = −z ± , (1.47)
√ 2
2
y(z) = ∓ . (1.48)
2
Note that in fact y − x − z = 0, so the Jacobian determinant ∂(f, g)/∂(x, y) = 0; thus, √ the above
expression for dy/dz is indeterminant. However, we see from the explicit expression y = ∓ 2/2 that
in fact, dy/dz = 0. The two original functions and their loci of intersection are plotted in Figure 1.1.
It is seen that the surface represented by the quadratic function is a open cylindrical tube, and that
represented by the linear function is a plane. Note that planes and cylinders may or may not intersect.
If they intersect, it is most likely that the intersection will be a closed arc. However, when the plane
is aligned with the axis of the cylinder, the intersection will be two non-intersecting lines; such is the
case in this example.
Let’s see how slightly altering the equation for the plane removes the degeneracy. Take now
5x + y + z = 0, (1.49)
2 2 2
x + y + z + 2xz = 1. (1.50)
Can x and y be considered as functions of z? If x = x(z) and y = y(z), then dx/dz and dy/dz must
exist. If we take
f (x, y, z) = 5x + y + z = 0, (1.51)
2 2 2
g(x, y, z) = x + y + z + 2xz − 1 = 0, (1.52)
T
then the solution matrix (dx/dz, dy/dz) is found as before:
∂f ∂f
− ∂z ∂y −1 1
∂g ∂g
dx − −(2z + 2x) 2y −2y + 2z + 2x
= ∂f∂z ∂f∂y = = , (1.53)
dz ∂x ∂y 5 1
10y − 2x − 2z
∂g ∂g 2x + 2z 2y
∂x ∂y
1 x
-0.2 0
1 0.2
y
0 0.5
y
0
-1 -0.5
-1
-2
1
1
0.5
z
0
z
0
-0.5
-1
-1
-0.5 -1
0
x 0.5
1
Figure 1.2: Surfaces of 5x+ y + z = 0 and x2 + y 2 + z 2 + 2xz = 1, and their loci of intersection
∂f
∂x − ∂f
∂z
5 −1
∂g ∂g 2x + 2z −(2z + 2x)
dy − ∂z −8x − 8z
= ∂x ∂f = = . (1.54)
dz ∂f 5 1 10y − 2x − 2z
∂x ∂y
∂g ∂g 2x + 2z 2y
∂x ∂y
The two original functions and their loci of intersection are plotted in Figure 1.2.
Straightforward algebra in this case shows that an explicit dependency exists:
√ √
−6z ± 2 13 − 8z 2
x(z) = , (1.55)
√26√
−4z ∓ 5 2 13 − 8z 2
y(z) = . (1.56)
26
These curves represent the projection of the curve of intersection on the x − z and y − z planes,
respectively. In both cases, the projections are ellipses.
2 2 2
(ds)2 = dξ 1 + dξ 2 + dξ 3 , (1.57)
3
X
2
(ds) = dξ idξ i ≡ dξ i dξ i. (1.58)
i=1
Here we have adopted the summation convention that when an index appears twice, a
summation from 1 to 3 is understood.
Now let us map a point from a point in (ξ 1, ξ 2 , ξ 3 ) space to a point in a more convenient
(x , x2 , x3 ) space. This mapping is achieved by defining the following functional dependen-
1
cies:
x1 = x1 (ξ 1, ξ 2 , ξ 3 ), (1.59)
x2 = x2 (ξ 1, ξ 2 , ξ 3 ), (1.60)
x3 = x3 (ξ 1 , ξ 2 , ξ 3). (1.61)
Taking derivatives can tell us whether the inverse exists.
∂x1 1 ∂x1 2 ∂x1 3 ∂x1 j
dx1 = dξ + 2 dξ + 3 dξ = dξ , (1.62)
∂ξ 1 ∂ξ ∂ξ ∂ξ j
2 2
∂x ∂x ∂x2 3 ∂x2 j
dx2 = dξ 1
+ dξ 2
+ dξ = dξ , (1.63)
∂ξ 1 ∂ξ 2 ∂ξ 3 ∂ξ j
∂x3 1 ∂x3 2 ∂x3 3 ∂x3 j
dx3 = dξ + 2 dξ + 3 dξ = dξ , (1.64)
∂ξ 1 ∂ξ ∂ξ ∂ξ j
1 ∂x1 ∂x1 ∂x1
dx ∂ξ 1 ∂ξ 2 ∂ξ 3 dξ 1
dx2 = ∂x2 ∂x2 ∂x2
∂ξ1 ∂ξ2 ∂ξ3 dξ 2 , (1.65)
dx3 ∂x3
1
∂x3
2
∂x3
3
dξ 3
∂ξ ∂ξ ∂ξ
i
∂x j
dxi = dξ . (1.66)
∂ξ j
In order for the inverse to exist we must have a non-zero Jacobian for the transformation,
i.e.
∂(x1 , x2 , x3 )
6= 0. (1.67)
∂(ξ 1 , ξ 2, ξ 3 )
4
Euclid of Alexandria, ∼ 325 B.C.-∼ 265 B.C., Greek geometer.
as one realizes the implications of the notation, however, the convention adopted ultimately does not matter.
6
Josiah Willard Gibbs, 1839-1903, prolific American physicist and mathematician with a lifetime affili-
ation with Yale University.
Example 1.4
Transform the Cartesian equation
∂S 2 2
1
− S = ξ1 + ξ2 . (1.86)
∂ξ
under the following:
1. Cartesian to linearly homogeneous affine coordinates.
Consider the following linear non-orthogonal transformation:
x1 = 2ξ 1 + ξ 2 , (1.87)
2
x = −8ξ 1 + ξ 2 , (1.88)
x3 = ξ3. (1.89)
2
ξ
4
1
x = constant
3
2
x = constant
2
0 1
0 1 2 3 4 ξ
Figure 1.3: Lines of constant x1 and x2 in the ξ 1 , ξ 2 plane for affine transformation of example
problem.
This transformation is of the class of affine transformations, which are of the form xi = Aij ξ j +bi . Affine
transformations for which bi = 0 are further distinguished as linear homogeneous transformations. The
transformation of this example is both affine and linear homogeneous.
This is a linear system of three equations in three unknowns; using standard techniques of linear
algebra allows us to solve for ξ 1 , ξ 2 , ξ 3 in terms of x1 , x2 , x3 ; that is we find the inverse transformation,
which is
1 1 1
ξ1 = x − x2 , (1.90)
10 10
4 1 1 2
ξ2 = x + x , (1.91)
5 5
ξ3 = x3 . (1.92)
Lines of constant x1 and x2 in the ξ 1 , ξ 2 plane are plotted in Figure 1.3. The appropriate Jacobian
matrix for the inverse transformation is
∂ξ 1 ∂ξ 1 ∂ξ 1
i 1 2 3 ∂x1 ∂x2 ∂x3
∂ξ ∂(ξ , ξ , ξ ) ∂ξ 2 ∂ξ 2 ∂ξ 2
J= j
= = , (1.93)
∂x ∂(x1 , x2 , x3 ) ∂x1
∂ξ 3
∂x2
∂ξ 3
∂x3
∂ξ 3
∂x1 ∂x2 ∂x
3
1 1
10 − 10 0
4 1
J = 5 5 0. (1.94)
0 0 1
So a unique transformation always exists, since the Jacobian determinant is never zero.
The metric tensor is
∂ξ i ∂ξ i ∂ξ 1 ∂ξ 1 ∂ξ 2 ∂ξ 2 ∂ξ 3 ∂ξ 3
gkl = = + + . (1.96)
∂xk ∂xl ∂xk ∂xl ∂xk ∂xl ∂xk ∂xl
∂ξ i ∂ξ i ∂ξ 1 ∂ξ 1 ∂ξ 2 ∂ξ 2 ∂ξ 3 ∂ξ 3
g11 = 1 1
= 1 1
+ 1 1+ 1 1 (1.97)
∂x ∂x
∂x ∂x ∂x ∂x ∂x ∂x
1 1 4 4 13
g11 = + + (0)(0) = . (1.98)
10 10 5 5 20
Repeating this operation for all terms of gkl , we find the complete metric tensor is
13 3
20 20 0
3 1
gkl = 20 20 0, (1.99)
0 0 1
13 1 3 3 1
g = det (gkl ) = − = . (1.100)
20 20 20 20 100
This is equivalent to the calculation in Gibbs notation:
G = JT · J (1.101)
1 4
1 1
10 5 0 10 − 10 0
G = − 1 1
0 · 4 1
0, (1.102)
10 5 5 5
0 0 1 0 0 1
13 3
20 20 0
3 1
G =
20 20 0. (1.103)
0 0 1
The details of the method of quadratic forms are delayed until a later chapter; direct expansion reveals
the two forms for (ds)2 to be identical. Note:
• The Jacobian matrix J is not symmetric.
• The metric tensor G = JT · J is symmetric.
• The fact that the metric tensor has non-zero off-diagonal elements is a consequence of the transfor-
mation being non-orthogonal.
• The distance is guaranteed to be positive. This will be true for all affine transformations in ordinary
three-dimensional Euclidean space. In the generalized space-time continuum suggested by the theory
of relativity, the generalized distance may in fact be negative; this generalized distance ds for an
2 2 2 2
infinitesimal change in space and time is given by ds2 = dξ 1 + dξ 2 + dξ 3 − dξ 4 , where the
2 2
first three coordinates are the ordinary Cartesian space coordinates and the fourth is dξ 4 = (c dt) ,
where c is the speed of light.
Note this system of equations is non-linear. For such systems, we cannot always find an explicit algebraic
expression for the inverse transformation. In this case, some straightforward algebraic and trigonometric
manipulation reveals that we can find an explicit representation of the inverse transformation, which is
ξ1 = x1 cos x2 , (1.112)
ξ2 = x1 sin x2 , (1.113)
ξ3 = x3 . (1.114)
Lines of constant x1 and x2 in the ξ 1 , ξ 2 plane are plotted in Figure 1.4. Notice that the lines of
constant x1 are orthogonal to lines of constant x2 in the Cartesian ξ 1 , ξ 2 plane. For general transfor-
mations, this will not be the case.
The appropriate Jacobian matrix for the inverse transformation is
x 2 = π/2
ξ2
x = 3π/4 3 x 2 = π/4
x1 = 3
2 x1 = 2
1 x1 = 1
x2 = π -2 -1 1 2 3 ξ
1
x 2= 0
-3
-1
-2
x 2 = 5π/4 -3
x 2 = 7π/4
2
x = 3π/2
Figure 1.4: Lines of constant x1 and x2 in the ξ 1 , ξ 2 plane for cylindrical transformation of
example problem.
∂ξ 1 ∂ξ 1 ∂ξ 1
i 1 2 3 ∂x1 ∂x2 ∂x3
∂ξ ∂(ξ , ξ , ξ ) ∂ξ 2 ∂ξ 2 ∂ξ 2
J= = = , (1.115)
∂xj ∂(x1 , x2 , x3 ) ∂x1
∂ξ 3
∂x2
∂ξ 3
∂x3
∂ξ 3
∂x1 ∂x2 ∂x3
cos x2 −x sin x2
1
0
J = sin x2 x1 cos x2 0. (1.116)
0 0 1
The determinant of the Jacobian matrix is
x1 cos2 x2 + x1 sin2 x2 = x1 . (1.117)
So a unique transformation fails to exist when x1 = 0.
The metric tensor is
∂ξ i ∂ξ i ∂ξ 1 ∂ξ 1 ∂ξ 2 ∂ξ 2 ∂ξ 3 ∂ξ 3
gkl = k l
= k l
+ k l + k l. (1.118)
∂x ∂x ∂x ∂x ∂x ∂x ∂x ∂x
For example for k = 1, l = 1 we get
∂ξ i ∂ξ i ∂ξ 1 ∂ξ 1 ∂ξ 2 ∂ξ 2 ∂ξ 3 ∂ξ 3
g11 = = + + , (1.119)
∂x1 ∂x1 ∂x1 ∂x1 ∂x1 ∂x1 ∂x1 ∂x1
g11 = cos2 x2 + sin2 x2 + 0 = 1. (1.120)
Repeating this operation, we find the complete metric tensor is
1 0 0
2
gkl = 0 x1 0, (1.121)
0 0 1
g = det (gkl ) = (x1 )2 . (1.122)
G = JT · J, (1.123)
cos x2 sin x2 0 cos x2 −x1 sin x2 0
G = −x sin x x cos x2
1 2 1
0 · sin x2 x1 cos x2 0, (1.124)
0 0 1 0 0 1
1 0 0
2
G = 0 x1 0. (1.125)
0 0 1
Note:
• The fact that the metric tensor is diagonal can be attributed to the transformation being orthogonal.
• Since the product of any matrix with its transpose is guaranteed to yield a symmetric matrix, the
metric tensor is always symmetric.
Also we have the volume ratio of differential elements as
Now
∂S ∂S ∂x1 ∂S ∂x2 ∂S ∂x3
= + + ,
∂ξ 1 ∂x1 ∂ξ 1 ∂x2 ∂ξ 1 ∂x3 ∂ξ 1
∂S ξ1 ∂S ξ2
= q − ,
∂x1 2 2 ∂x2 (ξ 1 )2 + (ξ 2 )2
(ξ 1 ) + (ξ 2 )
∂S sin x2 ∂S
= cos x2 − .
∂x1 x1 ∂x2
So the transformed version of Eq. (1.86) becomes
∂S sin x2 ∂S 2
cos x2 1
− 1 2
− S = x1 . (1.131)
∂x x ∂x
∂ x̄i j
ūi = u. (1.132)
∂xj
Quantities known as covariant vectors transform according to
∂xj
ūi = uj . (1.133)
∂ x̄i
Here we have considered general transformations from one non-Cartesian coordinate system
(x1 , x2 , x3 ) to another (x̄1 , x̄2 , x̄3 ).
Example 1.5
Let’s say (x, y, z) is a normal Cartesian system and define the transformation
Now we can assign velocities in both the unbarred and barred systems:
dx dy dz
ux = uy = uz =
dt dt dt
dx̄ dȳ dz̄
ūx̄ = ūȳ = ūz̄ =
dt dt dt
∂ x̄ dx ∂ ȳ dy ∂ z̄ dz
ūx̄ = ūȳ = ūz̄ =
∂x dt ∂y dt ∂z dt
ūx̄ = λux ūȳ = λuy ūz̄ = λuz
∂ x̄ x ∂ ȳ y ∂ z̄ z
ūx̄ = u ūȳ = u ūz̄ = u
∂x ∂y ∂z
This suggests the velocity vector is contravariant.
Now consider a vector which is the gradient of a function f (x, y, z). For example, let
f (x, y, z) = x + y 2 + z 3
∂f ∂f ∂f
ux = uy = uz =
∂x ∂y ∂z
ux = 1 uy = 2y uz = 3z 2
In the new coordinates x̄ ȳ z̄ x̄ ȳ 2 z̄ 3
f , , = + 2+ 3
λ λ λ λ λ λ
so
x̄ ȳ 2 z̄ 3
f¯ (x̄, ȳ, z̄) = + 2 + 3
λ λ λ
Now
∂ f¯ ∂ f¯ ∂ f¯
ūx̄ = ūȳ = ūz̄ =
∂ x̄ ∂ ȳ ∂ z̄
1 2ȳ 3z̄ 2
ūx̄ = ūȳ = ūz̄ =
λ λ2 λ3
In terms of x, y, z, we have
1 2y 3z 2
ūx̄ = ūȳ = ūz̄ =
λ λ λ
So it is clear here that, in contrast to the velocity vector,
1 1 1
ūx̄ = ux ūȳ = uy ūz̄ = uz
λ λ λ
Somewhat more generally we find for this case that
∂x ∂y ∂z
ūx̄ = ux ūȳ = uy ūz̄ = uz ,
∂ x̄ ∂ ȳ ∂ z̄
which suggests the gradient vector is covariant.
∂ x̄i ∂ x̄j kl
v̄ ij = v
∂xk ∂xl
Covariant tensors transform according to
∂xk ∂xl
v̄ij = vkl
∂ x̄i ∂ x̄j
Mixed tensors transform according to
∂ x̄i ∂xl k
v̄ji = v
∂xk ∂ x̄j l
Recall that variance is another term for gradient and that co- denotes with. A vector
which is co-variant is aligned with the variance or the gradient. Recalling next that contra-
denotes against, a vector with is contra-variant is aligned against the variance or the gra-
dient. This results in a set of contravariant basis vectors being tangent to lines of xi = C,
while covariant basis vectors are normal to lines of xi = C. A vector in space has two natural
representations, one on a contravariant basis, and the other on a covariant basis. The con-
travariant representation seems more natural, though both can be used to obtain equivalent
results. For the transformation x1 = (ξ 1 )2 + (ξ 2), x2 = (ξ 1) − (ξ 2 )3 , Figure 1.5 gives a sketch
of a set of lines of constant x1 and x2 in the Cartesian ξ 1 , ξ 2 plane, along with a local set of
both contravariant and covariant basis vectors.
The idea of covariant and contravariant derivatives play an important role in mathemat-
ical physics, namely in that the equations should be formulated such that they are invariant
under coordinate transformations. This is not particularly difficult for Cartesian systems,
but for non-orthogonal systems, one cannot use differentiation in the ordinary sense but
must instead use the notion of covariant and contravariant derivatives, depending on the
2
ξ
1
2
x = 1/16
0.8
2
covariant x = 1/2
basis vectors
0.6
contravariant
basis vectors
0.4
0.2 1
x =1
1
x = 1/2 1
ξ
0.2 0.4 0.6 0.8 1
Figure 1.5: Contours for the transformation x1 = (ξ 1 )2 + (ξ 2), x2 = (ξ 1 ) − (ξ 2)3 along with
a pair of contravariant basis vectors, which are tangent to the contours, and covariant basis
vectors, which are normal to the contours.
problem. The role of these terms was especially important in the development of the theory
of relativity.
Consider a contravariant vector ui defined in xi which has corresponding components U i
in the Cartesian ξ i . Take wji and Wji to represent the covariant spatial derivative of ui and
U i , respectively. Let’s use the chain rule and definitions of tensorial quantities to arrive at
a formula for covariant differentiation. From the definition of contravariance
∂ξ i l
Ui = u (1.135)
∂xl
Take the derivative in Cartesian space and then use the chain rule:
∂U i ∂U i ∂xk
Wji = = (1.136)
∂ξ j ∂xk ∂ξ j
i k
∂ ∂ξ l
= k
u ∂x (1.137)
∂x ∂xl ∂ξ j
| {z }
=U i
2 i
∂ ξ l ∂ξ i ∂ul ∂xk
= u + l k (1.138)
∂xk ∂xl ∂x ∂x ∂ξ j
2 p
∂ ξ ∂ξ p ∂ul ∂xk
Wqp = u l
+ (1.139)
∂xk ∂xl ∂xl ∂xk ∂ξ q
From the definition of a mixed tensor
∂xi ∂ξ q
wji = Wqp p j (1.140)
∂ξ ∂x
2 p
∂ ξ l ∂ξ p ∂ul ∂xk ∂xi ∂ξ q
= u + l k (1.141)
∂xk ∂xl ∂x ∂x ∂ξ q ∂ξ p ∂xj
| {z }
=Wqp
∂ 2 ξ p ∂xi
Γijl = (1.146)
∂xj ∂xl ∂ξ p
and use the term ∆j to represent the covariant derivative. Thus, the covariant derivative of
a contravariant vector ui is as follows:
∂ui
∆j ui = wji = + Γijl ul (1.147)
∂xj
Example 1.6
Find ∇T · u in cylindrical coordinates.9 The transformations are
q
2 2
x1 = + (ξ 1 ) + (ξ 2 )
2
ξ
x2 = tan−1
ξ1
x3 = ξ3
ξ1 = x1 cos x2
ξ2 = x1 sin x2
ξ3 = x3
∂ 2ξp ∂xi l
Γiil ul = u
∂xi ∂xl ∂ξ p
∂ 2ξ1 ∂xi l ∂ 2 ξ 2 ∂xi l ∂ 2 ξ 3 ∂xi l
= u + u + u
∂xi ∂xl ∂ξ 1 ∂xi ∂xl ∂ξ 2 ∂xi ∂xl ∂ξ 3
∂ 2 ξ 1 ∂xi l ∂ 2 ξ 2 ∂xi l
= u + u
∂xi ∂xl ∂ξ 1 ∂xi ∂xl ∂ξ 2
∂ 2 ξ 1 ∂x1 l ∂ 2 ξ 1 ∂x2 l ∂ 2 ξ 1 ∂x3 l
= u + u + u
∂x1 ∂xl ∂ξ 1 ∂x2 ∂xl ∂ξ 1 ∂x3 ∂xl ∂ξ 1
∂ 2 ξ 2 ∂x1 ∂ 2 ξ 2 ∂x2 ∂ 2 ξ 2 ∂x3
+ 1 l 2 ul + 2 l 2 ul + 3 l 2 ul
∂x ∂x ∂ξ ∂x ∂x ∂ξ ∂x ∂x ∂ξ
∂
∂ξ1
In Cartesian coordinates, we take ∇ ≡ ∂ξ∂ 2 . This gives rise to the natural, albeit unconventional,
9
∂
∂ξ3
T ∂ ∂ ∂
notation ∇ = ∂ξ1 ∂ξ2 ∂ξ3 .
∂ 2 ξ 1 ∂x1 l ∂ 2 ξ 1 ∂x2
= 1 l 1
u + 2 l 1 ul
∂x ∂x ∂ξ ∂x ∂x ∂ξ
∂ 2 ξ 2 ∂x1 l ∂ 2 ξ 2 ∂x2
+ 1 l 2 u + 2 l 2 ul
∂x ∂x ∂ξ ∂x ∂x ∂ξ
∂ 2 ξ 1 ∂x1 1 ∂ 2 ξ 1 ∂x1 2 ∂ 2 ξ 1 ∂x1 3
= u + u + u
∂x1 ∂x1 ∂ξ 1 ∂x1 ∂x2 ∂ξ 1 ∂x1 ∂x3 ∂ξ 1
∂ 2 ξ 1 ∂x2 ∂ 2 ξ 1 ∂x2 ∂ 2 ξ 1 ∂x2
+ 2 1 1 u1 + 2 2 1 u2 + 2 3 1 u3
∂x ∂x ∂ξ ∂x ∂x ∂ξ ∂x ∂x ∂ξ
∂ 2 ξ 2 ∂x1 1 ∂ 2 ξ 2 ∂x1 2 ∂ 2 ξ 2 ∂x1
+ 1 1 2 u + 1 2 2 u + 1 3 2 u3
∂x ∂x ∂ξ ∂x ∂x ∂ξ ∂x ∂x ∂ξ
∂ 2 ξ 2 ∂x2 1 ∂ 2 ξ 2 ∂x2 2 ∂ 2 ξ 2 ∂x2
+ 2 1 2 u + 2 2 2 u + 2 3 2 u3
∂x ∂x ∂ξ ∂x ∂x ∂ξ ∂x ∂x ∂ξ
∂ 2 ξ 1 ∂x1 1 ∂ 2 ξ 1 ∂x1 2
= u + u
∂x1 ∂x1 ∂ξ 1 ∂x1 ∂x2 ∂ξ 1
∂ 2 ξ 1 ∂x2 ∂ 2 ξ 1 ∂x2
+ 2 1 1 u1 + 2 2 1 u2
∂x ∂x ∂ξ ∂x ∂x ∂ξ
2 2 1
∂ ξ ∂x ∂ 2 ξ 2 ∂x1
+ 1 1 2 u1 + 1 2 2 u2
∂x ∂x ∂ξ ∂x ∂x ∂ξ
2 2 2
∂ ξ ∂x ∂ 2 ξ 2 ∂x2
+ 2 1 2 u1 + 2 2 2 u2
∂x ∂x ∂ξ ∂x ∂x ∂ξ
T 1 ∂ dr 1 ∂ dθ ∂ dz
∇ ·u = r + r +
r ∂r dt r ∂θ dt ∂z dt
1 ∂ 1 ∂uθ ∂uz
∇T · u = (rur ) + +
r ∂r r ∂θ ∂z
Here we have also used the more traditional uθ = r dθ 1 2 1 3
dt = x u , along with ur = u , uz = u . For
practical purposes, this insures that ur , uθ , uz all have the same dimensions.
Example 1.7
Calculate the acceleration vector du
dt in cylindrical coordinates.
Start by expanding the total derivative as
du ∂u
= + uT · ∇u.
dt ∂t
Now, we take u to be a contravariant velocity vector and the gradient operation to be a covariant
derivative. Employ index notation to get
du ∂ui
= + uj ∆j ui
dt ∂t i
∂ui j ∂u i l
= +u + Γjl u .
∂t ∂xj
After an extended calculation similar to the previous example, one finds after expanding all terms that
1 1 1
∂u 3 ∂u1 2
∂t u1 ∂u
∂x1 + u2 ∂u
∂x 2 + u ∂x 3 −x1 u2
du 2 1 ∂u2 2 2
= ∂u ∂t
+ u ∂x1 + u2 ∂u2 + u3 ∂u3 + 2 u u
∂x ∂x
1 2 .
dt ∂u3 3 3 3
x1
∂t u1 ∂u
∂x1 + u2 ∂u 3 ∂u
∂x2 + u ∂x3
0
The last term is related to the well known Coriolis10 and centripetal acceleration terms. However, these
are not in the standard form to which most are accustomed. To arrive at that standard form, one must
return to a so-called physical representation. Here again take x1 = r, x2 = θ, and x3 = z. Also take
ur = dr 1 dθ 1 2 dz 3
dt = u , uθ = r dt = x u , uz = dt = u . Then the r acceleration equation becomes
Here the final term is the traditional centripetal acceleration. The θ acceleration is slightly more
complicated. First one writes
dr dθ
d dθ ∂ dθ dr ∂ dθ dθ ∂ dθ dz ∂ dθ
= + + + + 2 dt dt .
dt dt ∂t dt dt ∂r dt dt ∂θ dt dt ∂z dt r
duθ ∂ dθ dr ∂ dθ dθ ∂ dθ dz ∂ dθ dr dθ
= r +r +r +r +2 ,
dt ∂t dt dt ∂r dt dt ∂θ dt dt ∂z dt dt dt
10
Gaspard-Gustave Coriolis, 1792-1843, French mechanician.
∂ dθ dr ∂ dθ dθ r dθ
dt ∂ dθ dz ∂ dθ dr r dθ
dt
= r + r − + r + r +2 ,
∂t dt dt ∂r dt dt r ∂θ dt dt ∂z dt dt r
∂uθ ∂uθ uθ ∂uθ ∂uθ ur uθ
= + ur + + uz + .
∂t ∂r r ∂θ ∂z r }
| {z
Coriolis
The final term here is the Coriolis acceleration. The z acceleration then is easily seen to be
duz ∂uz ∂uz uθ ∂uz ∂uz
= + ur + + uz .
dt ∂t ∂r r ∂θ ∂z
We summarize some useful identities, all of which can be proved, as well as some other
common notation, as follows
∂ξ i ∂ξ i
gkl = (1.148)
∂xk ∂xl
g = det(gij ) (1.149)
gik gkj
= δij (1.150)
ui = gij uj (1.151)
ui = g ij uj (1.152)
uT · v = ui v i = ui vi = gij uj v i = g ij uj vi (1.153)
u×v = ǫijk gjm gkn um v n = ǫijk uj vk (1.154)
∂ 2 ξ p ∂xi 1 ip ∂gpj ∂gpk ∂gjk
Γijk = = g + − (1.155)
∂xj ∂xk ∂ξ p 2 ∂xk ∂xj ∂xp
∂ui
∇u = ∆j ui = ui,j = j + Γijl ul (1.156)
∂x
∂ui 1 ∂ √ i
∇T · u = ∆i ui = ui,i = i + Γiil ul = √ gu (1.157)
∂x g ∂xi
p
ijk ijk p ijk ∂u p l
∇×u = ǫ uk,j = ǫ gkp u,j = ǫ gkp + Γjl u (1.158)
∂xj
du ∂u ∂ui ∂ui
= + uT · ∇u = + uj j + Γijl ul uj (1.159)
dt ∂t ∂t ∂x
∂φ
∇φ = φ,i = i (1.160)
∂x
∂ ij ∂φ ∂φ
∇2 φ = T ij
∇ · ∇φ = g φ,ij = j g i
+ Γjjk g ik i (1.161)
∂x ∂x ∂x
1 ∂ √ ij ∂φ
= √ j
gg (1.162)
g ∂x ∂xi
ij ∂T ij
∇T = T,k = k
+ Γilk T lj + Γjlk T il (1.163)
∂x
CC BY-NC-ND. 28 March 2011, M. Sen, J. M. Powers.
36 CHAPTER 1. MULTI-VARIABLE CALCULUS
∂T ij 1 ∂ √
∇T · T = T,jij = j
+ Γ i
lj T lj
+ Γj
lj T il
= √ j
g T ij
+ Γijk T jk (1.164)
∂x g ∂x
i
1 ∂ √ ∂ξ
= √ j
g T kj k (1.165)
g ∂x ∂x
Example 1.8
f = x2 − y 2
Equating partial derivatives with respect to x and to y to zero, we get
∂f
= 2x = 0
∂x
∂f
= −2y = 0
∂y
This gives x = 0, y = 0. For these values we find that
2 0
D = −
0 −2
= 4
Since D > 0, the point (0,0) is a saddle point.
11
Ludwig Otto Hesse, 1811-1874, German mathematician, studied under Jacobi.
Here t is a dummy variable of integration. Leibniz’s12 rule tells us how to take derivatives
of functions in integral form:
Z b(x)
y(x) = f (x, t) dt (1.168)
a(x)
Z b(x)
dy(x) db(x) da(x) ∂f (x, t)
= f (x, b(x)) − f (x, a(x)) + dt (1.169)
dx dx dx a(x) ∂x
Inverting this arrangement in a special case, we note if
Z x
y(x) = y(xo ) + f (t) dt (1.170)
x0
then (1.171)
Z x
dy(x) dx dxo ∂f (t)
= f (x) − f (x0 ) + dt (1.172)
dx dx dx x0 ∂x
dy(x)
= f (x) (1.173)
dx
Note that the integral expression naturally includes the initial condition that when x = x0 ,
y = y(x0 ). This needs to be expressed separately for the differential version of the equation.
Example 1.9
dy
Find dx if
Z x2
y(x) = (x + 1)t2 dt (1.174)
x
In this case, but not all, we can achieve the same result from explicit formulation of y(x):
Z x2
y(x) = (x + 1) t2 dt (1.180)
x
x2
t3
= (x + 1) (1.181)
3 x
x6 x3
= (x + 1) − (1.182)
3 3
7 6 4
x x x x3
y(x) = + − − (1.183)
3 3 3 3
dy(x) 7x6 4x 3
= + 2x5 − − x2 (1.184)
dx 3 3
So the two methods give identical results.
is an extremum. If y(x) is the desired solution, let Y (x) = y(x) + ǫh(x), where h(x1 ) =
h(x2 ) = 0. Thus Y (x) also satisfies the boundary conditions; also Y ′ (x) = y ′(x) + ǫh′ (x).
We can write Z x2
I(ǫ) = f (x, Y, Y ′ ) dx
x1
dI
Taking dǫ
, utilizing Leibniz’s rule, we get
Z x2
dI ∂f ∂x ∂f ∂Y ∂f ∂Y ′
= + + dx
dǫ x1 ∂x ∂ǫ ∂Y ∂ǫ ∂Y ′ ∂ǫ
Evaluating, we find
Z x2
dI ∂f ∂f ∂f ′
= 0+ h(x) + h (x) dx
dǫ x1 ∂x ∂Y ∂Y ′
1. f = f (x, y)
The Euler equation is
∂f
=0 (1.192)
∂y
which is easily solved:
f (x, y) = A(x) (1.193)
which, knowing f , is then solved for y(x).
2. f = f (x, y ′)
The Euler equation is
d ∂f
=0 (1.194)
dx ∂y ′
which yields
∂f
=A (1.195)
∂y ′
f (x, y ′) = Ay ′ + B(x) (1.196)
Again, knowing f , the equation is solved for y ′ and then integrated to find y(x).
3. f = f (y, y ′)
The Euler equation is
∂f d ∂f ′
− (y, y ) = 0 (1.197)
∂y dx ∂y ′
2
∂f ∂ f dy ∂ 2 f dy ′
− + =0 (1.198)
∂y ∂y∂y ′ dx ∂y ′ ∂y ′ dx
∂f ∂ 2 f dy ∂ 2 f d2 y
− − =0 (1.199)
∂y ∂y∂y ′ dx ∂y ′ ∂y ′ dx2
Multiply by y ′ to get
′ ∂f ∂ 2 f dy ∂ 2 f d2 y
y − − =0 (1.200)
∂y ∂y∂y ′ dx ∂y ′ ∂y ′ dx2
∂f ′′
Add and subtract ∂y ′y to get
′ ∂f ∂ 2 f dy ∂ 2 f d2 y ∂f ∂f
y − ′
− ′ ′ 2 + ′ y ′′ − ′ y ′′ = 0 (1.201)
∂y ∂y∂y dx ∂y ∂y dx ∂y ∂y
Regroup to get
2
∂f ′ ∂f ′′ ′ ∂ f dy ∂ 2 f d2 y ∂f ′′
y + ′y − y + + ′y = 0 (1.202)
∂y ∂y ∂y∂y ′ dx ∂y ′ ∂y ′ dx2 ∂y
| {z } | {z }
=df /dx =d/dx(y ′ ∂f /∂y ′ )
∂f
f (y, y ′) − y ′ =K (1.204)
∂y ′
Example 1.10
Find the curve of minimum length between the points (x1 , y1 ) and (x2 , y2 ).
If y(x) is the curve, then y(x1 ) = y1 and y(x2 ) = y2 . The length of the curve is
Z x2 p
L= 1 + (y ′ )2 dx
x1
Example 1.11
Find the curve through the points (x1 , y1 ) and (x2 , y2 ), such that the surface area of the body of
revolution by rotating the curve around the x-axis is a minimum.
We wish to minimize Z x2 p
I= y 1 + (y ′ )2 dx
x1
y 2
. 3
y -2
2.5
2
. 2
curve with
1.5
endpoints at
1 (-1, 3.09), (2, 2.26) z
0
which minimizes
0.5 surface area of body
of revolution corresponding
x -2
-1 -0.5 0 0.5 1 1.5 2 surface of
revolution
-1
0
1
x 2
Figure 1.6: Body of revolution of minimum surface area for (x1 , y1 ) = (−1, 3.08616) and
(x2 , y2 ) = (2, 2.25525)
p
Here f (y, y ′ ) = y 1 + (y ′ )2 . So the Euler equation reduces to
∂f
f (y, y ′ ) − y ′ = A
∂y ′
p y′
y 1 + y ′2 − y ′ y p = A
1 + y ′2
p
y(1 + y ′2 ) − yy ′2
A 1 + y ′2 =
p
y = A 1 + y ′2
r
′ y 2
y = −1
A
x−B
y(x) = A cosh
A
This is a catenary. The constants A and B are determined from the boundary conditions y(x1 ) = y1
and y(x2 ) = y2 . In general this requires a trial and error solution of simultaneous algebraic equations.
If (x1 , y1 ) = (−1, 3.08616) and (x2 , y2 ) = (2, 2.25525), one finds solution of the resulting algebraic
equations gives A = 2, B = 1.
For these conditions, the curve y(x) along with the resulting body of revolution of minimum surface
area are plotted in Figure 1.6.
Define
f ∗ = f − λ1 g 1 − λ2 g 2 − . . . − λn g n (1.206)
where the λi (i = 1, 2, · · · , n) are unknown constants called Lagrange14 multipliers. To get
the extremum of f ∗ , we equate to zero its derivative with respect to x1 , x2 , . . . , xm . Thus we
have
∂f ∗
= 0, i = 1, . . . , m (1.207)
∂xi
gi = 0, i = 1, . . . , n (1.208)
which are (m+n) equations that can be solved for xi (i = 1, 2, . . . , m) and λi (i = 1, 2, . . . , n).
Example 1.12
Extremize f = x2 + y 2 subject to the constraint 5x2 − 6xy + 5y 2 = 8.
Let
f ∗ = x2 + y 2 − λ(5x2 − 6xy + 5y 2 − 8)
from which
∂f ∗
= 2x − 10λx + 6λy = 0
∂x
∂f ∗
= 2y + 6λx − 10λy = 0
∂y
g = 5x2 − 6xy + 5y 2 = 8
x = ±y
√ √ √ √
The last equation gives the extrema to be at (x, y) = ( 2, 2), (− 2, − 2), ( √12 , − √12 ), (− √12 , √12 ).
The first two sets give f = 4 (maximum) and the last two f = 1 (minimum). The function to be
maximized along with the constraint function and its image are plotted in Figure 1.7.
A similar technique can be used for the extremization of a functional with constraint.
We wish to find the function y(x), with x ∈ [x1 , x2 ], and y(x1 ) = y1 , y(x2 ) = y2 , such that
the integral Z x2
I= f (x, y, y ′) dx (1.209)
x1
g=0 (1.210)
14
Joseph-Louis Lagrange, 1736-1813, Italian-born French mathematician.
x
2 x y 1 -1
y -1 0
1 0 1
0 1 0
2
-1
-1
-2
8 constrained
4
function
6 3
f(x,y) f(x,y)
4 2
1
2
0
0
unconstrained constraint
function function
Figure 1.7: Unconstrained function f (x, y) along with constrained function and constraint
function (image of constrained function)
Define
I ∗ = I − λg (1.211)
and continue as before.
Example 1.13
Extremize I, where Z a p
I= y 1 + (y ′ )2 dx
0
That is find the maximum surface area of a body of revolution which has a constant length. Let
Z ap
g= 1 + (y ′ )2 dx − ℓ = 0
0
Then let Z a Z a
∗
p p
I = I − λg = ′ 2
y 1 + (y ) dx − λ 1 + (y ′ )2 dx + λℓ
0 0
Z a p
= (y − λ) 1 + (y ′ )2 dx + λℓ
0
Z a
p λℓ
= (y − λ) 1 + (y ′ )2 + dx
0 a
x
0.2 0.4 0.6 0.8 1
-0.05
0.2
-0.1
-0.15
-0.2 0
y
-0.25
-0.3 -0.2
y
0.2 0
0 0.25
z 0.5
-0.2 x
0.75
1
Figure 1.8: Curve of length ℓ = 5/4 with y(0) = y(1) = 0 whose surface area of corresponding
body of revolution (also shown) is maximum.
p λℓ
With f ∗ = (y − λ) 1 + (y ′ )2 + a , we have the Euler equation
∂f ∗ d ∂f ∗
− =0
∂y dx ∂y ′
λℓ
Integrating from an earlier developed relationship, Eq. (1.204), when f = f (y, y ′ ), and absorbing a
into a constant A, we have
p y′
(y − λ) 1 + (y ′ )2 − y ′ (y − λ) p =A
1 + (y ′ )2
from which
p
(y − λ)(1 + (y ′ )2 ) − (y ′ )2 (y − λ) = A 1 + (y ′ )2
p
(y − λ) 1 + (y ′ )2 − (y ′ )2 = A 1 + (y ′ )2
p
y − λ = A 1 + (y ′ )2
s 2
y−λ
y′ = −1
A
x−B
y = λ + A cosh
A
Here A, B, λ have to be numerically determined from the three conditions y(0) = y(a) = 0, g = 0.
If we take the case where a = 1, ℓ = 5/4, we find that A = 0.422752, B = 12 , λ = −0.754549. For
these values, the curve of interest, along with the surface of revolution, is plotted in Figure 1.8.
Problems
1. If
z 3 + zx + x4 y = 2y 3 ,
(b) evaluate
∂z ∂z
, ,
∂x y ∂y x
2. Determine the general curve y(x), with x ∈ [x1 , x2 ], of total R x length L with endpoints y(x1 ) = y1
and y(x2 ) = y2 fixed, for which the area under the curve, x12 y dx, is a maximum. Show that if
(x1 , y1 ) = (0, 0); (x2 , y2 ) = (1, 1); L = 3/2, that the curve which maximizes the area and satisfies all
constraints is the circle, (y + 0.254272)2 + (x − 1.2453)2 = (1.26920)2. Plot this curve. What is the
area? Verify that each constraint is satisfied. What function y(x) minimizes the area and satisfies all
constraints? Plot this curve. What is the area? Verify that each constraint is satisfied.
3. Show that if a ray of light is reflected from a mirror, the shortest distance of travel is when the angle
of incidence on the mirror is equal to the angle of reflection.
4. The speed of light in different media separated by a planar interface is c1 and c2 . Show that if the
time taken for light to go from a fixed point in one medium to another in the second is a minimum,
the angle of incidence, αi , and the angle of refraction, αr , are related by
sin αi c1
=
sin αr c2
5. F is a quadrilateral with perimeter P . Find the form of F such that its area is a maximum. What is
this area?
6. A body slides due to gravity from point A to point B along the curve y = f (x). There is no friction
and the initial velocity is zero. If points A and B are fixed, find f (x) for which the time taken will
be the least. What is this time? If A : (x, y) = (1, 2), B : (x, y) = (0, 0), where distances are in
meters, plot the minimum time curve, and find the minimum time if the gravitational acceleration is
g = −9.81 sm2 j.
R1
7. Consider the integral I = 0 (y ′ − y + ex )2 dx. What kind of extremum does this integral have
(maximum or minimum)? What should y(x) be for this extremum? What does the solution of the
Euler equation give, if y(0) = 0 and y(1) = −e? Find the value of the extremum. Plot y(x) for the
extremum. If y0 (x) is the solution of the Euler equation, compute I for y1 (x) = y0 (x) + h(x), where
you can take any h(x) you like, but with h(0) = h(1) = 0.
8. Find the length of the shortest curve between two points with cylindrical coordinates (r, θ, z) = (a, 0, 0)
and (r, θ, z) = (a, Θ, Z) along the surface of the cylinder r = a.
9. Determine the shape of a parallelogram with a given area which has the least perimeter.
with y(0) = 0 and y(1) = 1. Plot y(x) which renders the integral at an extreme point.
11. Find the point on the plane ax + by + cz = d which is nearest to the origin.
12. Extremize the integral Z 1
y ′2 dx
0
subject to the end conditions y(0) = 0, y(1) = 0, and also the constraint
Z 1
y dx = 1
0
Plot the function y(x) which extremizes the integral and satisfies all constraints.
13. Show that the functions
x+y
u =
x−y
xy
v =
(x − y)2
are functionally dependent.
14. Find the point on the curve of intersection of z − xy = 10 and x + y + z = 1, that is closest to the
origin.
15. Find a function y(x) with y(0) = 1, y(1) = 0 that extremizes the integral
r 2
Z 1 1 + dy
dx
I= dx
0 y
Plot y(x) for this function.
16. For elliptic cylindrical coordinates
ξ1 = cosh x1 cos x2
ξ2 = sinh x1 sin x2
3
ξ = x3
Find the Jacobian matrix J and the metric tensor G. Find the inverse transformation. Plot lines of
constant x1 and x2 in the ξ 1 and ξ 2 plane.
17. For the elliptic coordinate system of the previous problem, find ∇T · u where u is an arbitrary vector.
18. For parabolic coordinates
ξ1 = x1 x2 cos x3
ξ2 = x1 x2 sin x3
1
ξ3 = (x2 )2 − (x1 )2
2
Find the Jacobian matrix J and the metric tensor G. Find the inverse transformation. Plot lines of
constant x1 and x2 in the ξ 1 and ξ 2 plane.
19. For the parabolic coordinate system of the previous problem, find ∇T · u where u is an arbitrary
vector.
20. Find the covariant derivative of the contravariant velocity vector in cylindrical coordinates.
F (x, y, y ′) = 0 (2.1)
dy
where y ′ = dx
.
Example 2.1
Solve
8x + 1
yy ′ = , with y(1) = −5.
y
Separating variables
y 2 dy = 8xdx + dx.
Integrating, we have
y3
= 4x2 + x + C.
3
49
50 CHAPTER 2. FIRST-ORDER ORDINARY DIFFERENTIAL EQUATIONS
10
7.5
2.5
x
-10 -5 5 10
-2.5
-5
Figure 2.1: y(x) which solves yy ′ = (8x + 1)/y with y(1) = −5.
y 3 = 12x2 + 3x − 140.
Example 2.2
Solve
y2
xy ′ = 3y + , with y(1) = 4.
x
This can be written as y y 2
y′ = 3 + .
x x
Let u = y/x. Then
f (u) = 3u + u2
Using our developed formula, we get
du dx
2
= .
2u + u x
Since
1 1 1
= −
2u + u2 2u 4 + 2u
both sides can be integrated to give
1
(ln |u| − ln |2 + u|) = ln |x| + C.
2
1
The initial condition gives C = 2 ln 32 , so that the solution can be reduced to
y 2 2
2x + y = 3 x .
This can be solved explicitly for y(x) for each case of the absolute value. The first case
4 3
3x
y(x) =
1 − 23 x2
is seen to satisfy the condition at x = 1. The second case is discarded as it does not satisfy the condition
at x = 1.
The solution is plotted in Figure 2.2.
y
20
15
10
x
-6 -4 -2 2 4 6
-5
-10
-15
-20
y2
Figure 2.2: y(x) which solves xy ′ = 3y + x
with y(1) = 4
dF (x, y) = 0, (2.10)
where F (x, y) = 0 is a solution to the differential equation. The chain rule is used to expand
the derivative of F (x, y) as
∂F ∂F
dF = dx + dy = 0
∂x ∂y
So, for an equation of the form
As long as F (x, y) is continuous and differentiable, the mixed second partials are equal, thus,
∂P ∂Q
= (2.14)
∂y ∂x
must hold if F (x, y) is to exist and render the original differential equation to be exact.
Example 2.3
Solve
dy ex−y
=
dx ex−y −1
ex−y dx + 1 − ex−y dy = 0
| {z } | {z }
=P =Q
∂P
= −ex−y
∂y
∂Q
= −ex−y
∂x
∂P ∂Q
Since ∂y = ∂x , the equation is exact. Thus
∂F
= P (x, y)
∂x
∂F
= ex−y
∂x
F (x, y) = ex−y + A(y)
∂F dA
= −ex−y + = Q(x, y) = 1 − ex−y
∂y dy
dA
= 1
dy
A(y) = y−C
x−y
F (x, y) = e +y−C = 0
ex−y + y = C
Example 2.4
Solve
dy 2xy
= 2
dx x − y2
Separating variables, we get
(x2 − y 2 ) dy = 2xy dx.
y
6
C=2
2
C=1
C=0 x
-6 -4 -2 2 4 6
C = -1
C = -2 -2
This is not exact according to criterion (2.14). It turns out that the integrating factor is y −2 , so that
on multiplication, we get 2
2x x
dx − − 1 dy = 0.
y y2
This can be written as
x2
d +y =0
y
which gives
x2
+ y = C,
y
x2 + y 2 = Cy.
The solution for various values of C is plotted in Figure 2.4.
y
3
C=3
2 C=2
1 C=1
x
-1.5 -1 -0.5 0.5 1 1.5
-1 C = -1
C = -2
-2
C = -3
-3
2xy
Figure 2.4: y(x) which solves y ′(x) = (x2 −y 2 )
which yields
Rx
R xo
Z x Rt
− P (s)ds P (s)ds P (s)ds
y(x) = e a e a yo + ea Q(t)dt . (2.21)
xo
Example 2.5
Solve
y ′ − y = e2x ; y(0) = yo .
Here
P (x) = −1
or
P (s) = −1
yo = -2
yo = 0
yo = 2
y
3
x
-3 -2 -1 1 2 3
-1
-2
-3
dy(x)
e−x − e−x y(x) = ex
dx
d
e−x y(x) = ex
dx
d −t
e y(t) = et
Z x dt Z x
d −t
e y(t) dt = et dt
xo =0 dt xo =0
e−x y(x) − e−0 y(0) = e −e x 0
−x
e y(x) − yo = ex − 1
y(x) = ex (yo + ex − 1)
y(x) = e2x + (yo − 1) ex
Since S(x) is itself a solution to equation (2.24), we subtract appropriate terms to get
1 dz 2S 1 1
− 2 = P + 2 +Q (2.29)
z dx z z z
dz
− = P (2Sz + 1) + Qz (2.30)
dx
dz
+ (2P (x)S(x) + Q(x)) z = −P (x). (2.31)
dx
Again this is a first order linear equation in z and x of the form of equation (2.15) and can
be solved.
Example 2.6
Solve
e−3x 2 1
y′ = y − y + 3e3x .
x x
One solution is
y = S(x) = e3x ,
Verify:
e−3x 6x 1 3x
3e3x = e − e + 3e3x
x x
e3x e3x
3e3x = − + 3e3x
x x
3e3x = 3e3x
so let
1
y = e3x + .
z
Also we have
e−3x
P (x) =
x
1
Q(x) = −
x
R(x) = 3e3x
Substituting in the equation, we get
−3x
dz e 3x 1 e−3x
+ 2 e − z = −
dx x x x
dz z e−3x
+ =− .
dx x x
The integrating factor here is
dx
R
e x = eln x = x
Multiplying by the integrating factor x
dz
x + z = −e−3x
dx
C= -2
C= 0
C= 2
C= -2 C= -1 y
3
2.5
1.5
0.5
x
-1 -0.8 -0.6 -0.4 -0.2 0.2 0.4
-0.5
-1
C= -2 C= -1
exp(−3x)
Figure 2.6: y(x) which solves y ′ = x
− y/x + 3 exp(3x)
d(xz)
= −e−3x
dx
which can be integrated as
e−3x C e−3x + 3C
z= + = .
3x x 3x
Since y = S(x) + z1 , the solution is thus
3x
y = e3x + .
e−3x + 3C
The solution for various values of C is plotted in Figure 2.6.
2.7.1 y absent
If
f (x, y ′ , y ′′) = 0 (2.32)
then let u(x) = y ′ . Thus u′ (x) = y ′′ , and the equation reduces to
du
f x, u, =0 (2.33)
dx
Example 2.7
Solve
xy ′′ + 2y ′ = 4x3 .
Let u = y ′ , so that
du
x + 2u = 4x3 .
dx
Multiplying by x
du
x2 + 2xu = 4x4
dx
d 2
(x u) = 4x4 .
dx
This can be integrated to give
4 3 C1
u= x + 2
5 x
from which
1 4 C1
y= x − + C2
5 x
for x 6= 0.
2.7.2 x absent
If
f (y, y ′, y ′′ ) = 0 (2.34)
let u(x) = y ′, so that
dy ′ dy ′ dy du
y ′′ = = = u
dx dy dx dy
The equation becomes
du
f y, u, u =0 (2.35)
dy
which is also an equation of first order. Note however that the independent variable is now
y while the dependent variable is u.
Example 2.8
Solve
y ′′ − 2yy ′ = 0; y(0) = yo , y ′ (0) = yo′ .
du dy du
Let u = y ′ , so that y ′′ = dx = dx dy = u du
dy . The equation becomes
du
u − 2yu = 0
dy
Now
u=0
satisfies the equation. Thus
dy
= 0
dx
y = C
applying one initial condition: y = yo
This satisfies the initial conditions only under special circumstances, i.e. yo′ = 0. For u 6= 0,
du
= 2y
dy
u = y 2 + C1
apply I.C.’s: yo′ = yo2 + C1
C1 = yo′ − yo2
dy
= y 2 + yo′ − yo2
dx
dy
= dx
y 2 + yo′ − yo2
from which for yo′ − yo2 > 0
!
1 y
p tan−1 p = x + C2
yo − yo2
′ yo − yo2
′
!
1 −1 yo
p tan p = C2
yo′ − yo2 yo′ − yo2
!!
p p yo
y(x) = yo′ − yo2 tan x yo′ − yo2 + tan−1 p
yo′ − yo2
For yo′ − yo2 < 0, one would obtain solutions in terms of hyperbolic trigonometric functions.
y
3
x
-1.5 -1 -0.5 0.5 1 1.5
-1
-2
-3
y = eCx , (2.36)
y ′ = CeCx , (2.37)
y eCx
ln y = ln eCx , (2.38)
x x
eCx
= Cx, (2.39)
x
= CeCx , (2.40)
= y′. (2.41)
So the differential equation is satisfied for all values of C. Now to satisfy the initial condition,
we must have
2 = eC(0) , (2.42)
2 = 1? (2.43)
There is no finite value of C that allows satisfaction of the initial condition. The original
differential equation can be written as xy ′ = y ln y. The point x = 0 is singular since at that
point, the highest derivative is multiplied by 0 leaving only 0 = y ln y at x = 0. For the very
special initial condition y(0) = 1, the solution y = eCx is valid for all values of C. Thus for
this singular equation, for most initial conditions, no solution exists. For one special initial
condition, a solution exists, but it is not unique.
Theorem
Let f (x, y) be continuous and satisfy |f (x, y)| ≤ m and the Lipschitz condition |f (x, y) −
f (x, y0 )| ≤ k|y − y0 | in a bounded region R. Then the equation y ′ = f (x, y) has one and
only one solution containing the point (x0 , y0).
A stronger condition is that if f (x, y) and ∂f /∂y are finite and continuous at (x0 , y0 ),
then a solution of y ′ = f (x, y) exists and is unique in the neighborhood of this point.
Example 2.9
Analyze the uniqueness of the solution of
dy √
= −K y, y(T ) = 0.
dt
Here, t is the independent variable instead of x. Taking,
√
f (t, y) = −K y
we have
∂f K
=− √
∂y 2 y
which is not finite at y = 0. So the solution cannot be guaranteed to be unique. In fact, one solution is
1 2
y(t) = K (t − T )2 .
4
Another solution which satisfies the initial condition and differential equation is
y(t) = 0.
Example 2.10
Consider the equation
dy
= 3y 2/3 , with y(2) = 0.
dx
On separating variables and integrating
3y 1/3 = 3x + 3C
y
1
0.75
0.5
0.25
x
1 2 3 4
-0.25
-0.5
-0.75
-1
y
1
0.75
0.5
0.25
x
1 2 3 4
-0.25
-0.5
-0.75
-1
Figure 2.8: Two solutions y(x) which satisfy y ′ = 3y 2/3 with y(2) = 0
y = xy ′ + f (y ′) (2.44)
3
Alexis Claude Clairaut, 1713-1765, Parisian/French mathematician.
y = xu + f (u). (2.45)
There are two possible solutions to this, u′ = 0 or x + df /du = 0. If we consider the first
and take
du
u′ = = 0, (2.49)
dx
we can integrate to get
u=C (2.50)
where C is a constant. Then, from equation (2.45), we get the general solution
y = Cx + f (C) (2.51)
Applying an initial condition y(xo ) = yo gives what we will call the regular solution.
But if we take
df
x+ =0 (2.52)
du
then this equation along with equation (2.45)
df
y = −u + f (u) (2.53)
du
form a set of parametric equations for what we call the singular solution. It is singular
because the coefficient on the highest derivative in Eq. (2.48) is itself 0.
Example 2.11
Solve
y = xy ′ + (y ′ )3 , y(0) = yo
Take
u = y′
then
f (u) = u3
df
= 3u2
du
y
6 yo = 3
yo = 2
4
yo = 0
(singular) yo = 1
2
x yo = 0
-4 -3 -2 -1 1 2
-2
yo = 0
y o= -1
(singular)
-4
y o= -2
-6 y o= -3
Figure 2.9: Two solutions y(x) which satisfy y = xy ′ + (y ′)3 with y(0) = yo
Problems
1. Find the general solution of the differential equation
y ′ + x2 y(1 + y) = 1 + x3 (1 + x).
4. Solve
dy x−y
= .
dx x+y
5. Solve the nonlinear equation (y ′ − x)y ′′ + 2y ′ = 2x.
6. Solve xy ′′ + 2y ′ = x. Plot a solution for y(1) = 1, y ′ (1) = 1.
7. Solve y ′′ − 2yy ′ = 0. Plot a solution for y(0) = 0, y ′ (0) = 3.
8. Given that y1 = x−1 is one solution of y ′′ + x3 y ′ + 1
x2 y = 0, find the other solution.
9. Solve
(a) y ′ tan y + 2 sin x sin( π2 + x) + ln x = 0
(b) xy ′ − 2y − x4 − y 2 = 0
(c) y ′ cos y cos x + sin y sin x = 0
(d) y ′ + y cot x = ex
2
(e) x5 y ′ + y + ex (x6 − 1)y 3 = 0, with y(1) = e−1/2
(f) y ′ + y 2 − xy − 1 = 0
(g) y ′ (x + y 2 ) − y = 0
x+2y−5
(h) y ′ = −2x−y+4
(i) y ′ + xy = y
Plot solutions, when possible, for y(0) = −1, 0, 1.
10. Find all solutions of
(x + 1)(y ′ )2 + (x − y)y ′ − y = 0
where y(x) is an unknown function. The equation is said to be homogeneous if f (x) = 0, giving then
L(y) = 0 (3.2)
d d2
The operator L is composed of a combination of derivatives dx , dx2 etc. L is linear if
and
L(αy) = αL(y) (3.4)
where α is a scalar. The general form of L is
dn dn−1 d
L = Pn (x) n
+ Pn−1 (x) n−1 + . . . + P1 (x) + P0 (x) (3.5)
dx dx dx
The ordinary differential equation (3.1) is then linear.
Definition: The functions y1 (x), y2 (x), . . . , yn (x) are said to be linearly independent when C1 y1 (x)+C2 y2 (x)+
. . . + Cn yn (x) = 0 is true only when C1 = C2 = . . . = Cn = 0.
A homogeneous equation of order n can be shown to have n linearly independent solutions. These are
called complementary functions. If yi (i = 1, . . . , n) are the complementary functions of the equation, then
n
X
y(x) = Ci yi (x) (3.6)
i=1
69
70 CHAPTER 3. LINEAR ORDINARY DIFFERENTIAL EQUATIONS
is the general solution of the homogeneous equation. If yp (x) is a particular solution of equation (3.1), the
general solution is then
Xn
y(x) = yp (x) + Ci yi (x). (3.7)
i=1
Now we would like to show that any solution φ(x) to the homogeneous equation L(y) = 0 can be written
as a linear combination of the n complementary functions yi (x):
For a unique solution, we need the determinant of the coefficient matrix to be non-zero. This particular
determinant is known as the Wronskian1 W of y1 (x), y2 (x), . . . , yn (x) and is defined as
y1 y2 ... yn
y1′ y2′ ... yn′
W = .. .. .. (3.13)
(n−1). . ... .
y (n−1) (n−1)
1 y2 . . . yn
W 6= 0 indicates linear independence of the functions y1 (x), y2 (x), . . . , yn (x), since if φ(x) ≡ 0, the only
solution is Ci = 0, i = 1, . . . , n. Unfortunately, the converse is not always true; that is if W = 0, the
complementary functions may or may not be linearly dependent, though in most cases W = 0 indeed implies
linear dependence.
Example 3.1
Determine the linear independence of (a) y1 = x and y2 = 2x, (b) y1 = x and y2 = x2 , and (c)
y1 = x2 and y2 = x|x| for x ∈ (−1, 1)
x 2x
(a) W = = 0, linearly dependent.
1 2
x x2
(b) W = = x2 6= 0, linearly independent, except at x = 0.
1 2x
(c) We can restate y2 as
y2 (x) = −x2 x ∈ (−1, 0]
1
Józef Maria Hoene-Wroński, 1778-1853, Polish-born French mathematician.
y2 (x) = x2 x ∈ (0, 1)
so that 2
x −x2
W = = −2x3 + 2x3 = 0 x ∈ (−1, 0]
2x −2x
2
x x2
W = = 2x3 − 2x3 = 0 x ∈ (0, 1)
2x 2x
Thus W = 0 for x ∈ (−1, 1), which suggests the functions may be linearly dependent. However, when
we seek C1 and C2 such that C1 y1 + C2 y2 = 0, we find the only solution is C1 = 0, C2 = 0; therefore,
the functions are in fact linearly independent, despite the fact that W = 0! Let’s check this. For
x ∈ (−1, 0],
C1 x2 + C2 (−x2 ) = 0,
so we will need C1 = C2 at a minimum. For x ∈ (0, 1),
C1 x2 + C2 x2 = 0,
which gives the requirement that C1 = −C2 . Substituting the first condition into the second gives
C2 = −C2 , which is only satisfied if C2 = 0, thus requiring that C1 = 0; hence, the functions are indeed
linearly independent.
where Ai , (i = 0, . . . , n) are constants. To find the solution of this equation we let y = erx .
Substituting we get
This is called the characteristic equation. It is an nth order polynomial which has n roots
(some of which could be repeated, some of which could be complex), ri (i = 1, . . . , n)
from which n linearly independent complementary functions yi(x) (i = 1, . . . , n) have to be
obtained. The general solution is then given by equation (3.6).
If all roots are real and distinct, then the complementary functions are simply eri x ,
(i = 1, . . . , n). If, however, k of these roots are repeated, i.e. r1 = r2 = . . . = rk = r,
then the linearly independent complementary functions are obtained by multiplying erx by
1, x, x2 , . . . , xk−1 . For a pair of complex conjugate roots p ± qi, one can use de Moivre’s for-
mula (see Appendix) to show that the complementary functions are epx cos qx and epx sin qx.
Example 3.2
Solve
d4 y d3 y d2 y dy
4
− 2 3
+ 2
+2 − 2y = 0
dx dx dx dx
Substituting y = erx , we get a characteristic equation
r4 − 2r3 + r2 + 2r − 2 = 0
ay ′ + by = 0 (3.18)
is
ar + b = 0 (3.19)
so
b
r=− (3.20)
a
thus the complementary function for this equation is simply
b
y = Ce− a x (3.21)
d2 y dy
a 2
+ b + cy = 0 (3.22)
dx dx
is
ar 2 + br + c = 0 (3.23)
Depending on the coefficients of this quadratic equation, there are three cases to be consid-
ered.
• b2 − 4ac > 0: two distinct real roots r1 and r2 . The complementary functions are
y1 = er1 x and y2 = er2 x .
• b2 − 4ac = 0: one real root. The complementary functions are y1 = erx and y2 = xerx .
• b2 − 4ac < 0: two complex conjugate roots p ± qi. The complementary functions are
y1 = epx cos qx and y2 = epx sin qx.
Example 3.3
Solve
d2 y dy
2
−3 + 2y = 0
dx dx
The characteristic equation is
r2 − 3r + 2 = 0
with solutions
r1 = 1, r2 = 2.
The general solution is then
y = C1 ex + C2 e2x
Example 3.4
Solve
d2 y dy
2
−2 +y =0
dx dx
The characteristic equation is
r2 − 2r + 1 = 0
with repeated roots
r1 = 1, r2 = 1.
The general solution is then
y = C1 ex + C2 xex
Example 3.5
Solve
d2 y dy
2
−2 + 10y = 0
dx dx
The characteristic equation is
r2 − 2r + 10 = 0
with solutions
r1 = 1 + 3i, r2 = 1 − 3i.
The general solution is then
y = ex (C1 cos 3x + C2 sin 3x)
let the other solution be y2 (x) = u(x)y1 (x). We then form derivatives of y2 and substitute
into the original differential equation. First compute the derivatives:
which is solved for v(x) using known methods for first order equations.
z = ln x
so that
x = ez .
Then
dz 1
= = e−z ,
dx x
dy dy dz dy d d
= = e−z so = e−z
dx dz dx dz dx dz
2
2
dy d dy −z d −z dy −2z d y dy
= =e e =e −
dx2 dx dx dz dz dz 2 dz
Substituting into the differential equation, we get
d2 y dy
2
+ (A − 1) + By = 0 (3.33)
dz dz
which is an equation with constant coefficients.
In what amounts to the same approach, one can alternatively assume a solution of the
form y = Cxr . This leads to a characteristic polynomial for r of
r(r − 1) + Ar + B = 0. (3.34)
The two roots for r induce two linearly independent complementary functions.
Example 3.6
Solve
x2 y ′′ − 2xy ′ + 2y = 0, for x > 0.
With x = ez , we get
d2 y dy
2
− 3 + 2y = 0.
dz dz
The solution is
y = C1 ez + C2 e2z = C1 x + C2 x2 .
Note that this equation can also be solved by letting y = Cxr . Substituting into the equation, we get
r2 − 3r + 2 = 0, so that r1 = 1 and r2 = 2. The solution is then obtained as a linear combination of
xr1 and xr2 .
Example 3.7
Solve
d2 y dy
x2 + 3x + 15y = 0.
dx2 dx
Let us assume here that y = Cxr . Substituting this assumption into the equation yields
r(r − 1) + 3r + 15 = 0.
r2 + 2r + 15 = 0.
Solving gives √
r = −1 ± i 14.
Thus, we see there are two linearly independent complementary functions:
√ √
y(x) = C1 x−1+i 14
+ C2 x−1−i 14
.
Factoring gives
1 √ √
y(x) = C1 xi 14 + C2 x−i 14 .
x
Expanding in terms of exponentials and logarithms gives
1 √ √
y(x) = C1 (exp(ln x))i 14 + C2 (exp(ln x))−i 14 .
x
1 √ √
y(x) = C1 exp(i 14 ln x) + C2 exp(i 14 ln x) .
x
1 √ √
y(x) = Ĉ1 cos( 14 ln x) + Ĉ2 sin( 14 ln x) .
x
Example 3.8
y ′′ + 4y ′ + 4y = 169 sin 3x
Thus
r2 + 4r + 4 = 0
(r + 2)(r + 2) = 0
r1 = −2, r2 = −2
Since the roots are repeated, the complementary functions are
y1 = e−2x y2 = xe−2x
yp = a sin 3x + b cos 3x
so
yp′ = 3a cos 3x − 3b sin 3x
(−9a sin 3x − 9b cos 3x) + 4 (3a cos 3x − 3b sin 3x) + 4 (a sin 3x + b cos 3x) = 169 sin 3x
Example 3.9
Solve
y ′′′′ − 2y ′′′ + y ′′ + 2y ′ − 2y = x2 + x + 1
Let the particular integral be of the form yp = ax2 + bx + c. Substituting we get
−(2a + 1)x2 + (4a − 2b − 1)x + (2a + 2b − 2c − 1) = 0
For this to hold for all values of x, the coefficients must be zero, from which a = − 12 , b = − 23 , and
c = − 25 . Thus
1
yp = − (x2 + 3x + 5)
2
The solution of the homogeneous equation was found in a previous example, so that the general solution
is
1
y = C1 e−x + C2 ex + ex (C3 cos x + C4 sin x) − (x2 + 3x + 5)
2
Example 3.10
Solve
y ′′ + 4y = 6 sin 2x
Since sin 2x is a complementary function, we will try
yp = x(a sin 2x + b cos 2x)
from which
yp′ = 2x(a cos 2x − b sin 2x) + (a sin 2x + b cos 2x)
yp′′ = −4x(a sin 2x + b cos 2x) + 4(a cos 2x − b sin 2x)
Substituting into the equation, we compare coefficients and get a = 0, b = − 32 . The general solution
is then
3
y = C1 sin 2x + C2 cos 2x − x cos 2x.
2
Example 3.11
Solve
y ′′ + 2y ′ + y = xe−x
The complementary functions are e−x and xe−x . To get the particular solution we have to choose
a function of the kind yp = ax3 e−x . On substitution we find that a = 1/6. Thus the general solution is
1
y = C1 e−x + C2 xe−x + x3 e−x
6
we propose
n
X
yp = ui (x)yi (x) (3.36)
i=1
where yi(x), (i = 1, . . . , n) are complementary functions of the equation, and ui (x) are n
unknown functions. Differentiating, we have
n
X n
X
yp′ = u′i yi + ui yi′ .
i=1 i=1
Pn
We set i=1 u′i yi to zero as a first condition. Differentiating the rest
n
X n
X
yp′′ = u′i yi′ + uiyi′′ .
i=1 i=1
Again we set the first term on the right side to zero as a second condition. Following this
procedure repeatedly we arrive at
n
X n
X
(n−2) (n−1)
yp(n−1) = u′iyi + ui yi .
i=1 i=1
The vanishing of the first term on the right gives us the (n − 1)’th condition. Substituting
these in the governing equation, the last condition
n
X n
X
(n−1) (n) (n−1)
Pn (x) u′i yi + ui Pn y i + Pn−1 yi + ...+ P1 yi′ + P0 yi = f (x)
i=1 i=1 | {z }
=0
is obtained. Since each of the functions yi is a complementary function, the term within
brackets is zero.
To summarize, we have the following n equations in the n unknowns u′i , (i = 1, . . . , n)
that we have obtained:
n
X
u′iyi = 0,
i=1
n
X
u′iyi′ = 0,
i=1
..
. (3.37)
n
X (n−2)
u′i yi = 0,
i=1
n
X (n−1)
Pn (x) u′i yi = f (x).
i=1
These can be solved for u′i , and then integrated to give the ui ’s.
Example 3.12
Solve
y ′′ + y = tan x.
The complementary functions are
y1 = cos x, y2 = sin x
u′1 y1 + u′2 y2 = 0
u′1 y1′ + u′2 y2′ = tan x.
Integrating, we get
Z
u1 = − sin x tan x dx = sin x − ln | sec x + tan x|,
Z
u2 = cos x tan x dx = − cos x.
yp = u 1 y1 + u 2 y2
= (sin x − ln | sec x + tan x|) cos x − cos x sin x
= − cos x ln | sec x + tan x|
Ly = f (x), (3.38)
where the highest derivative in L is order n and with general homogeneous boundary condi-
tions at x = a and x = b on linear combinations of y and n − 1 of its derivatives:
T T
A y(a), y ′(a), . . . , y (n−1) (a) + B y(b), y ′(b), . . . , y (n−1) (b) =0 (3.39)
where A and B are n × n constant coefficient matrices, then knowing L, A and B, we can
form a solution of the form:
Z b
y(x) = f (s)g(x, s)ds (3.40)
a
This is desirable as
• once g(x, s) is known, the solution is defined for all f including
• numerical solution of the quadrature problem is more robust than direct numerical
solution of the original differential equation
• the solution is useful in experiments in which the system dynamics are well charac-
terized (e.g. mass spring damper) but the forcing may be erratic (perhaps digitally
specified)
We now define the Green’s2 function: g(x, s) and proceed to show that with this definition,
we are guaranteed to achieve the solution to the differential equation in the desired form as
shown at the beginning of the section. We take g(x, s) to be the Green’s function for the
linear differential operator L, as defined by Eq. (3.5), if it satisfies the following conditions:
1. Lg(x, s) = δ(x − s)
4. g(x, s), g ′(x, s), . . . , g (n−2) (x, s) are continuous for x ∈ [a, b]
5. g (n−1) (x, s) is continuous for [a, b] except at x = s where it has a jump of 1/Pn (s); the
jump is defined from left to right.
Also for purposes of the above conditions, s is thought of as a constant parameter. In the
actual Green’s function representation of the solution, s is a dummy variable. The Dirac
delta function δ(x − s) is discussed in the appendix and in Sec. 7.20 in Kaplan.
These conditions are not all independent; nor is the dependence obvious. Consider for
example,
d2 d
L = P2 (x) 2 + P1 (x) + Po (x)
dx dx
Then we have
d2 g dg
P2 (x) 2
+ P1 (x) + Po (x)g = δ(x − s)
dx dx
d2 g P1 (x) dg Po (x) δ(x − s)
+ + g =
dx2 P2 (x) dx P2 (x) P2 (x)
Integrating
Z
dg dg P1 (s) Po (s) s+ǫ 1 s+ǫ
− + g|s+ǫ − g|s−ǫ + g dx = H(x − s)|s−ǫ
dx s+ǫ dx s−ǫ P2 (s)
P2 (s) s−ǫ P2 (s)
This is consistent with the final point, that the second highest derivative of g suffers a jump
at x = s.
Next, we show that applying this definition of g(x, s) to our desired result lets us recover
the original differential equation, rendering g(x, s) to be appropriately defined. This can be
The analysis can be extended in a straightforward manner to more arbitrary systems with
inhomogeneous boundary conditions using matrix methods (c.f. Wylie and Barrett, 1995).
Example 3.13
Find the Green’s function and the corresponding solution integral of the differential equation
d2 y
= f (x)
dx2
subject to boundary conditions
y(0) = 0, y(1) = 0
Verify the solution integral if f (x) = 6x.
Here
d2
L=
dx2
Now 1) break the problem up into two domains: a) x < s, b) x > s, 2) Solve Lg = 0 in both domains;
four constants arise, 3) Use boundary conditions for two constants, 4) use conditions at x = s: continuity
dg
of g and a jump of dx , for the other two constants.
a) x < s
d2 g
= 0
dx2
dg
= C1
dx
g = C1 x + C2
g(0) = 0 = C1 (0) + C2
C2 = 0
g(x, s) = C1 x, x<s
b) x > s
d2 g
= 0
dx2
dg
= C3
dx
g = C3 x + C4
g(1) = 0 = C3 (1) + C4
C4 = −C3
g(x, s) = C3 (x − 1) , x>s
Continuity of g(x, s) when x = s:
C1 s = C3 (s − 1)
s−1
C1 = C3
s
s−1
g(x, s) = C3 x, x<s
s
g(x, s) = C3 (x − 1) , x>s
dg
Jump in dx at x = s (note P2 (x) = 1):
dg dg
− = 1
dx s+ǫ dx s−ǫ
s−1
C3 − C3 = 1
s
C3 = s
g(x, s) = x(s − 1), x<s
g(x, s) = s(x − 1), x>s
Note some properties of g(x, s) which are common in such problems:
• it is broken into two domains
• it is continuous in and through both domains
• its n − 1 (here n = 2, so first) derivative is discontinuous at x = s
• it is symmetric in s and x across the two domains
• it is seen by inspection to satisfy both boundary conditions
The general solution in integral form can be written by breaking the integral into two pieces as
Z x Z 1
y(x) = f (s) s(x − 1) ds + f (s) x(s − 1) ds
0 x
Z x Z 1
y(x) = (x − 1) f (s) s ds + x f (s) (s − 1) ds
0 x
Now evaluate the integral if f (x) = 6x (thus f (s) = 6s).
Z x Z 1
y(x) = (x − 1) (6s) s ds + x (6s) (s − 1) ds
0 x
Z x Z 1
2
= (x − 1) 6s ds + x 6s2 − 6s ds
0 x
x 1
= (x − 1) 2s3 0 + x 2s3 − 3s2 x
= (x − 1)(2x3 − 0) + x((2 − 3) − (2x3 − 3x2 ))
= 2x4 − 2x3 − x − 2x4 + 3x3
y(x) = x3 − x
Note the original differential equation and both boundary conditions are automatically satisfied by the
solution.
The solution is plotted in Figure 3.1.
-0.1
0.5
-0.2 x
-2 -1 1 2
-0.5
-0.3
-1
-1.5
3
y(x) = x - x 3
y(x) = x - x
in domain of interest 0 < x < 1 in expanded domain, -2 < x < 2
3.3.4 Operator D
The linear operator D is defined by
dy
D(y) = .
dx
or, in terms of the operator alone,
d
D=
dx
The operator can be repeatedly applied, so that
dn y
Dn (y) = .
dxn
Another example of its use is
dy(x)
= f (x) y(xo ) = yo
dx Z x
y(x) = yo + f (s) ds
xo
then
This gives us h(x) explicitly in terms of the known function f such that h satisfies D(h)−ah =
f.
We can iterate to find the solution to higher order equations such as
Note that Z x
dy
= yobeb(x−xo ) + h(x) + bebx h(s)e−bs ds
dx xo
dy
(xo ) = yo′ = yo b + h(xo )
dx
which can be rewritten as
(D − b)(y(xo )) = h(xo )
which is what one would expect.
Returning to the problem at hand, we take our expression for h(x), evaluate it at x = s
and substitute into the expression for y(x) to get
Z x Z s
b(x−xo ) bx a(s−xo ) as −at
y(x) = yo e +e h(xo )e +e f (t)e dt e−bs ds
xo xo
Z x Z s
b(x−xo ) bx ′ a(s−xo ) as
= yo e +e (yo − yo b) e +e f (t)e dt e−bs ds
−at
xo xo
Z x Z s
b(x−xo ) bx ′ (a−b)s−axo (a−b)s −at
= yo e +e (yo − yo b) e +e f (t)e dt ds
xo xo
Z x Z x Z s
b(x−xo ) bx ′ (a−b)s−axo bx (a−b)s −at
= yo e + e (yo − yo b) e ds + e e f (t)e dt ds
xo xo xo
Z x Z s
b(x−xo ) bx ′ ea(x−xo )−xb − e−bxo bx (a−b)s −at
= yo e + e (yo − yo b) +e e f (t)e dt ds
a−b xo xo
Z x Z s
b(x−xo ) ′ ea(x−xo ) − eb(x−xo ) bx (a−b)s −at
= yo e + (yo − yob) +e e f (t)e dt ds
a−b xo xo
Z xZ s
b(x−xo ) ea(x−xo ) − eb(x−xo ) bx
= yo e ′
+ (yo − yob) +e e(a−b)s f (t)e−at dt ds
a−b xo xo
Z xZ x
b(x−xo ) ea(x−xo ) − eb(x−xo ) bx
= yo e + (yo′
− yob) +e e(a−b)s f (t)e−at ds dt
a−b xo t
a(x−xo ) b(x−xo ) Z x Z x
b(x−xo ) ′ e −e bx −at (a−b)s
= yo e + (yo − yob) +e f (t)e e ds dt
a−b xo t
Z x
b(x−xo ) ea(x−xo ) − eb(x−xo ) f (t) a(x−t)
= yo e ′
+ (yo − yob) + e − eb(x−t) dt
a−b xo a − b
Thus we have a solution to the second order linear differential equation with constant
coefficients and arbitrary forcing expressed in integral form. A similar alternate expression
can be developed when a = b.
Problems
1. Find the general solution of the differential equation
y ′ + x2 y(1 + y) = 1 + x3 (1 + x)
2. Show that the functions y1 = sin x, y2 = x cos x, and y3 = x are linearly independent. Find the lowest
order differential equation of which they are the complementary functions.
3. Solve the following initial value problem for (a) C = 6, (b) C = 4, and (c) C = 3 with y(0) = 1 and
y ′ (0) = −3.
d2 y dy
+C + 4y = 0
dt2 dt
Plot your results.
4. Solve
d3 y 2
d y
(a) dx3 − 3 dx 2 + 4y = 0
d4 y 3
d y 2
d y dy
(b) dx4 − 5 dx 3 + 11 dx2 − 7 dx = 12
8. Solve
x2 y ′′ + xy ′ − 4y = 6x
y ′′ + y ′ − 2y = f (x)
with y(0) = 0, y ′ (1) = 0. Determine y(x) if f (x) = 3 sin x. Plot your result.
11. Find the Green’s function solution of
y ′′ + 4y = f (x)
with y(0) = y(1), y ′ (0) = 0. Verify this is the correct solution when f (x) = x2 . Plot your result.
12. Solve y ′′′ − 2y ′′ − y ′ + 2y = sin2 x.
13. Solve y ′′′ + 6y ′′ + 12y ′ + 8y = ex − 3 sin x − 8e−2x .
14. Solve x4 y ′′′′ + 7x3 y ′′′ + 8x2 y ′′ = 4x−3 .
15. Show that x−1 and x5 are solutions of the equation
x2 y ′′ − 3xy ′ − 5y = 0
This chapter will deal with series solution methods. Such methods are useful in solving both
algebraic and differential equations. The first method is formally exact in that an infinite
number of terms can often be shown to have absolute and uniform convergence properties.
The second method, asymptotic series solutions, is less rigorous in that convergence is not
always guaranteed; in fact convergence is rarely examined because the problems tend to
be intractable. Still asymptotic methods will be seen to be quite useful in interpreting the
results of highly non-linear equations in local domains.
91
92 CHAPTER 4. SERIES SOLUTION METHODS
Example 4.1
Find the power series solution of
dy
=y y(0) = yo
dx
around x = 0. Let
y = a0 + a1 x + a2 x2 + a3 x3 + · · ·
so that
dy
= a1 + 2a2 x + 3a3 x2 + 4a4 x3 + · · ·
dx
Substituting in the equation, we have
If this is valid for all x, the coefficients must be all zero. Thus
a1 = a0
1 1
a2 = a1 = a0
2 2
1 1
a3 = a2 = a0
3 3!
1 1
a4 = a3 = a0
4 4!
..
.
so that
x2 x3 x4
y(x) = a0 1 + x + + + + ···
2! 3! 4!
Applying the initial condition at x = 0 gives ao = yo so
x2 x3 x4
y(x) = yo 1 + x + + + + ···
2! 3! 4!
Of course this power series is the Taylor1 series expansion of the closed form solution y = yo ex .
For yo = 1 the exact solution and three approximations to the exact solution are shown in Figure
4.1.
1
Brook Taylor, 1685-1731, English mathematician, musician, and painter.
y
y’ = y y = exp( x)
4
y (0) = 1 y = 1 + x + x 2/ 2
3
y=1+x
2
1 y=1
x
-1.5 -1 -0.5 0.5 1 1.5
2
Karl Theodor Wilhelm Weierstrass, 1815-1897, Westphalia-born German mathematician.
such that
|un (x)| ≤ Mn
for all x in the domain. For our problem, we take the domain to be −A ≤ x ≤ A, where A > 0.
So for uniform convergence we must have
n
x
≤ Mn
n!
So take
An
Mn =
n!
(Note Mn is thus strictly positive). So
∞ ∞
X X An
Mn =
n=0 n=0
n!
A
lim ≤1
n→∞ n + 1
This holds for all A, so in the domain, −∞ < x < ∞ the series converges absolutely and uniformly.
d2 y dy
P (x) 2 + Q(x) + R(x)y = 0 (4.3)
dx dx
around x = a. There are three different cases, depending of the behavior of P (a), Q(a) and
R(a), in which x = a is classified as a ordinary point, a regular singular point, or an irregular
singular point. These are described below.
Example 4.2
Find the series solution of
y ′′ + xy ′ + y = 0 y(0) = yo y ′ (0) = yo′
around x = 0.
x = 0 is an ordinary point, so that we have
∞
X
y = an xn
n=0
X∞
y′ = nan xn−1
n=1
X∞
xy ′ = nan xn
n=1
X∞
xy ′ = nan xn
n=0
X∞
y ′′ = n(n − 1)an xn−2
n=2
X∞
m=n−2 = (m + 1)(m + 2)am+2 xm
m=0
∞
X
= (n + 1)(n + 2)an+2 xn
n=0
-1
y = 1 - x 2 /2
The series converges for all x. For yo = 1, yo′ = 0 the exact solution, which can be shown to be
2
x
y = exp − ,
2
and two approximations to the exact solution are shown in Figure 4.2.
3
Ferdinand Georg Frobenius, 1849-1917, Prussian/German mathematician.
2. r1 = r2 = r. Then
∞
X
r
y1 = (x − a) an (x − a)n (4.6)
n=0
∞
X
r
y2 = y1 ln x + (x − a) bn (x − a)n (4.7)
n=1
The constants an and k are determined by the differential equation. The general solution is
y(x) = C1 y1 (x) + C2 y2 (x) (4.10)
Example 4.3
Find the series solution of
4xy ′′ + 2y ′ + y = 0
around x = 0.
x = 0 is a regular singular point. So we take
∞
X
y = an xn+r
n=0
X∞
y′ = an (n + r)xn+r−1
n=0
X∞
y ′′ = an (n + r)(n + r − 1)xn+r−2
n=0
∞
X ∞
X ∞
X
4 an (n + r)(n + r − 1)xn+r−1 + 2 an (n + r)xn+r−1 + an xn+r = 0
n=0 n=0 n=0
| {z } | {z } | {z }
=4xy ′′ =2y ′ =y
∞
X ∞
X
2 an (n + r)(2n + 2r − 1)xn+r−1 + an xn+r = 0
n=0 n=0
∞
X X∞
m=n−1 2 am+1 (m + 1 + r)(2(m + 1) + 2r − 1)xm+r + an xn+r = 0
m=−1 n=0
X∞ X∞
2 an+1 (n + 1 + r)(2(n + 1) + 2r − 1)xn+r + an xn+r = 0
n=−1 n=0
r(2r − 1) = 0
For r = 0
1
an+1 = −an
(2n + 2)(2n + 1)
x x2 x3
y1 = a0 1 − + − + ···
2! 4! 6!
1
For r = 2
1
an+1 = −an
2(2n + 3)(n + 1)
1/2 x x2 x3
y 2 = a0 x 1− + − + ···
3! 5! 7!
√ √
The series converges for all x to y1 = cos x and y2 = sin x. The general solution is
y = C1 y1 + C2 y2
or √ √
y(x) = C1 cos x + C2 sin x
Note that y(x) is real and non-singular for x ∈ [0, ∞). However, the first derivative
√ √
′ sin x cos x
y (x) = −C1 √ + C2 √
2 x 2 x
is singular at x = 0. The nature of the singularity is seen from a Taylor series expansion of y ′ (x) about
x = 0, which gives
√
1 x 1 x
y ′ (x) ∼ C1 − + + . . . + C2 √ − + ...
2 12 2 x 4
√
So there is a weak 1/ x singularity in y ′ (x) at x = 0.
′
For y(0) = 1, y (0) < ∞, the exact solution and the linear approximation to the exact solution are
shown in Figure 4.3. For this case, one has C1 = 1 to satisfy the condition on y(0), and one must have
C2 = 0 to satisfy the non-singular condition on y ′ (0).
Example 4.4
Find the series solution of
xy ′′ − y = 0
y
y = cos (x1/2 ) (exact)
1
x
20 40 60 80 100
-1
-2 4 x y’’ + 2 y’ + y = 0
-3 y (0) = 1
y=1-x/2
y ’ (0) <
8
-4
around x = 0.
P∞
Let y = n=0 an xn+r . Then, from the equation
∞
X
r−1
r(r − 1)a0 x + ((n + r)(n + r − 1)an − an−1 ) xn+r−1 = 0
n=1
1
an = an−1
n(n + 1)
1
= a0
n!(n + 1)!
Thus
1 1 1 4
y1 (x) = x + x2 + x3 + x + ...
2 12 144
The second solution is
∞
X
y2 (x) = ky1 (x) ln x + bn xn
n=0
k = b0
1 k(2n + 1)
bn+1 = bn − for n = 1, 2, . . .
n(n + 1) n!(n + 1)!
Thus
3 2 7 3 35 4 1 2 1 3 1 4
y2 (x) = b0 y1 ln x + b0 1 − x − x − x − . . . + b1 x + x + x + x + ...
4 36 1728 2 12 144
| {z }
=y1
Since the last series is y1 (x), we choose b0 = 1 and b1 = 0. The general solution is
1 1 1 4
y(x) = C1 x + x2 + x3 + x + ...
2 12 144
1 2 1 3 1 4 3 7 35 4
+C2 x+ x + x + x + . . . ln x + 1 − x2 − x3 − x − ...
2 12 144 4 36 1728
where I1 and K1 are what is known as modified Bessel functions of the first and second kinds, respec-
tively, both of order 1. The function I1 (s) is non-singular, while K1 (s) is singular at s = 0.
Example 4.5
Solve
y ′′′ − xy = 0
around x = 0.
Let
∞
X
y= an xn
n=0
exact y = 1 + x 4 / 24
7
6
5
y’’’ - x y = 0,
4 y(0) = 1,
3 y’ (0) = 0,
y’’ (0) = 0.
2
1
x
-4 -2 2 4
from which
∞
X
xy = an−1 xn
n=1
∞
X
y ′′′ = 6a3 + (n + 1)(n + 2)(n + 3)an+3 xn
n=1
a3 = 0
1
an+3 = an−1
(n + 1)(n + 2)(n + 3)
which gives the general solution
1 4 1 8
y(x) = a0 1 + x + x + ...
24 8064
1 4 1 8
+a1 x 1 + x + x + ...
60 30240
2 1 4 1 8
+a2 x 1 + x + x + ...
120 86400
For yo = 1, y ′ (0) = 0, y ′′ (0) = 0 the exact solution and the linear approximation to the exact solution are
shown in Figure 4.4. The exact solution is expressed in terms of a generalized hypergeometric function
whose form is non-standard. The software package Mathematica gives more details, and reports the
exact solution as 2
1 3 x
y = HypergeometricPFQ {} , , , .
2 4 64
Example 4.6
For 0 < ǫ << 1 solve
x2 + ǫx − 1 = 0
Let
x = x0 + ǫx1 + ǫ2 x2 + · · ·
Substituting in the equation,
2
x0 + ǫx1 + ǫ2 x2 + · · · +ǫ x0 + ǫx1 + ǫ2 x2 + · · · −1 = 0
| {z } | {z }
=x2 =x
x
exact 2
3 x +εx-1=0
linear
2
-3 -2 -1 1 2 3 ε
-1
-2
linear
-3
exact
Example 4.7
For 0 < ǫ << 1 solve
ǫx2 + x − 1 = 0
Note as ǫ → 0, the equation becomes singular. Let
x = x0 + ǫx1 + ǫ2 x2 + · · ·
O(ǫ0 ) : x0 − 1 = 0 ⇒ x0 = 1
O(ǫ1 ) : x20 + x1 = 0 ⇒ x1 = −1
O(ǫ2 ) : 2x0 x1 + x2 = 0 ⇒ x2 = 2
..
.
X = xǫ X2 + X − ǫ = 0
We expand
X = X0 + ǫX1 + ǫ2 X2 + · · ·
so 2
X0 + ǫX1 + ǫ2 X2 + · · · + X0 + ǫX1 + ǫ2 X2 + · · · − ǫ = 0
X02 + 2ǫX0 X1 + ǫ2 (X12 + 2X0 X2 ) + · · · + X0 + ǫX1 + ǫ2 X2 + · · · − ǫ = 0
Collecting terms of the same order
X = −1 − ǫ + ǫ2 + · · ·
X = ǫ − ǫ2 + · · ·
or
1
x = −1 − ǫ + ǫ2 + · · ·
ǫ
x = 1 − ǫ + ···
Example 4.8
Solve
cos x = ǫ sin(x + ǫ)
π
for x near 2.
Figure 4.7 shows a plot of cos x and ǫ sin(x + ǫ) for ǫ = 0.1. It is seen that there are multiple
intersections near x = n + 21 π . We seek only one of these.
x
3
1
asymptotic exact
-1 1 2 3
ε
-1
exact
asymptotic
-2 asymptotic
-3
f(x)
ε = 0.1 1
cos (x)
......
0.5
ε sin(x + ε)
x
-10 -5 5 10
-0.5
-1
We substitute
x = x0 + ǫx1 + ǫ2 x2 + · · ·
we have
cos(x0 + ǫx1 + ǫ2 x2 + · · ·) = ǫ sin(x0 + ǫx1 + ǫ2 x2 + · · · + ǫ)
Now we expand both the left and right hand sides in a Taylor series in ǫ about ǫ = 0. We note that
a general function f (ǫ) has such a Taylor series of f (ǫ) ∼ f (0) + ǫf ′ (0) + (ǫ2 /2)f ′′ (0) + . . . Expanding
the left hand side, we get
= d/dǫ(cos x)|ǫ=0
z }| {
cos(x0 + ǫx1 + . . .) = cos(x0 + ǫx1 + . . .)|ǫ=0 +ǫ (− sin(x0 + ǫx1 + . . .)) (x1 + 2ǫx2 + . . .) +...,
| {z } | {z } | {z } | {z }
=cos x = cos x| =d/dx(cos x)| = dx/dǫ|
ǫ=0 ǫ=0 ǫ=0 ǫ=0
Collecting terms
O(ǫ0 ) : cos x0 = 0 ⇒ x0 = π2
O(ǫ1 ) : −x1 sin x0 − sin x0 = 0 ⇒ x1 = −1
..
.
The solution is
π
x= − ǫ + ···
2
Example 4.9
For 0 < ǫ << 1 solve
y ′′ + ǫy 2 = 0, with y(0) = 1, y ′ (0) = 0
Let
y(x) = y0 (x) + ǫy1 (x) + ǫ2 y2 (x) + · · ·
y ′ (x) = y0′ (x) + ǫy1′ (x) + ǫ2 y2′ (x) + · · ·
y ′′ (x) = y0′′ (x) + ǫy1′′ (x) + ǫ2 y2′′ (x) + · · ·
Substituting in the equation
2
y0′′ (x) + ǫy1′′ (x) + ǫ2 y2′′ (x) + · · · + ǫ y0 (x) + ǫy1 (x) + ǫ2 y2 (x) + · · · = 0
y0′′ (x) + ǫy1′′ (x) + ǫ2 y2′′ (x) + · · · + ǫ y02 (x) + 2ǫy1 (x)yo (x) + · · · = 0
The solution is
x2 x4
y =1−ǫ + ǫ2 + ···
2 12
For validity of the asymptotic solution, we must have
x2
1 >> ǫ .
2
This solution becomes invalid when the first term is as large or larger than the second:
x2
, 1≤ǫ
2
r
2
|x| ≥ .
ǫ
Using the techniques of the previous chapter it is seen that this equation has an exact solution.
With
dy d2 y dy ′ dy du
u= 2
= = u
dx dx dy dx dy
the original equation becomes
du
u + ǫy 2 = 0
dy
udu = −ǫy 2 dy
u2 ǫ
= − y 3 + C1
2 3
ǫ
u=0 when y=1 so C=
3
r
2ǫ
u=± (1 − y 3 )
3
r
dy 2ǫ
=± (1 − y 3 )
dx 3
dy
dx = ± q
2ǫ 3
3 (1 − y )
Z y
ds
x=± q
2ǫ 3
3 (1 − s )
1
asymptotic x
-6 -4 -2 2 4 6
y’’ + ε y2 = 0 -10
y(0) = 1 -0.5
y’(0) = 0
-15
ε = 0.1 exact -1
asymptotic
-20 -1.5
It can be shown that this integral can be represented in terms of Gauss’s4 hypergeometric function,
2 F1 (a, b, c, z) as follows:
r r
π Γ 13 3 1 1 4 3
x=∓ ± y 2 F1 , , ,y
6ǫ Γ 56 2ǫ 3 2 3
It is likely difficult to invert either of the above functions to get y(x) explicitly. For small ǫ, the essence
of the solution is better conveyed by the asymptotic solution. A portion of the asymptotic and exact
solutions for ǫ = p 0.1 are shown in Figure 4.8. For this value, the asymptotic solution is expected to be
invalid for |x| ≥ 2/ǫ = 4.47.
Example 4.10
Solve
y ′′ + ǫy 2 = 0, with y(0) = 1, y ′ (0) = ǫ
Let
y(x) = y0 (x) + ǫy1 (x) + ǫ2 y2 (x) + · · ·
Substituting in the equation and collecting terms
O(ǫ0 ) : y0′′ = 0, y0 (0) = 1, y0′ (0) = 0 ⇒ y0 = 1
2
O(ǫ1 ) : y1′′ = −y02 , y1 (0) = 0, y1′ (0) = 1 ⇒ y1 = − x2 + x
4 3
O(ǫ2 ) : y2′′ = −2y0 y1 , y2 (0) = 0, y2′ (0) = 0 ⇒ y2 = x12 − x3
..
.
The solution is 4
x2 x x3
− x + ǫ2
y =1−ǫ − + ···
2 12 3
A portion of the asymptotic and exact solutions for ǫ = 0.1 are shown in Figure 4.9. Compared to the
previous example, there is a slight offset from the y axis.
4
Johann Carl Friedrich Gauss, 1777-1855, Brunswick-born German mathematician of tremendous influ-
ence.
y
1
x
-10 -5 5 10
-1
-2
-3
y’’ + ε y 2 = 0 exact
y(0) = 1
-4 y’(0) = ε
ε = 0.1
-5
asymptotic
Example 4.11
Find an approximate solution of the Duffing equation:
First let’s give some physical motivation, as also outlined in Section 10.2 of Kaplan. One problem in
which Duffing’s equation arises is the undamped motion of a mass subject to a non-linear spring force.
Consider a body of mass m moving in the horizontal x plane. Initially the body is given a small positive
displacement x(0) = xo . The body has zero initial velocity dx dt (0) = 0. The body is subjected to a
non-linear spring force Fs oriented such that it will pull the body towards x = 0:
Fs = (ko + k1 x2 )x
Here ko and k1 are dimensional constants with SI units N/m and N/m3 respectively. Newton’s second
law gives us
d2 x
m 2 = −(ko + k1 x2 )x
dt
d2 x dx
m + (ko + k1 x2 )x = 0, x(0) = xo , (0) = 0
dt2 dt
Choose an as yet arbitrary length scale L and an as yet arbitrary time scale T with which to scale the
problem and take:
x t
x̃ = t̃ =
L T
Substitute
mL d2 x̃ L dx̃
+ ko Lx̃ + k1 L3 x̃3 = 0 Lx̃(0) = xo , (0) = 0
T 2 dt̃2 T dt̃
x
6 3
x’’ + x + ε x = 0
4 x(0) = 1, x’(0) = 0
t
20 40 60 80 100
-2
-4
ε = 0.2
-6
d2 x̃ ko T 2 k1 L2 T 2 3 xo dx̃
2
+ x̃ + x̃ = 0 x̃(0) = , (0) = 0
dt̃ m m L dt̃
Now we want to examine the effect of small non-linearities. Choose the length and time scales such
that the leading order motion has an amplitude which is O(1) and a frequency which is O(1). So take
r
m
T ≡ L ≡ xo
ko
So
d2 x̃ k1 x2o kmo 3 dx̃
+ x̃ + x̃ = 0 x̃(0) = 1, (0) = 0
dt̃2 m dt̃
Choosing
k1 x2o
ǫ≡
ko
we get
d2 x̃ dx̃
+ x̃ + ǫx̃3 = 0 x̃(0) = 1, (0) = 0
dt̃2 dt̃
So our asymptotic theory will be valid for
Now, let’s drop the superscripts and focus on the mathematics. First, a very accurate numerical
approximation to the exact solution x(t) for ǫ = 0.2 is shown Figure 4.10. The so-called phase plane for
this solution, giving x versus dx/dt is shown in Figure 4.11. Note if ǫ = 0, the solution is x(t) = cos t,
and thus dx/dt = − sin t. Thus for ǫ = 0, x2 + (dx/dt)2 = cos2 t+ sin2 t = 1. Thus the ǫ = 0 phase plane
solution is a unit circle. Figure 4.11 displays a small deviation from a circle. This deviation would be
more pronounced for larger ǫ.
Let’s use an asymptotic method to try to capture this solution. Using the expansion
dx
dt
1.0
0.5
x
-1.0 -0.5 0.5 1.0
-0.5
-1.0
Figure 4.11: Numerical solution in the phase plane, dx/dt versus x, to Duffing’s equation,
ǫ = 0.2.
The difference between the exact solution and the leading order solution, xexact (t) − xo (t) is plotted in
Figure 4.12. The error is the same order of magnitude as the solution itself for moderate values of t.
This is undesirable.
To O(ǫ) the solution is
ǫ
x = cos t + (− cos t + cos 3t − 12t sin t) + · · ·
32
The original differential equation can be integrated once via the following steps
ẋ ẍ + x + ǫx3 = 0,
ẋẍ + ẋx + ǫẋx3 = 0,
d 1 2 1 2 ǫ 4
ẋ + x + x = 0,
dt 2 2 4
1 2 1 2 ǫ 4 1 2 1 2 ǫ 4
ẋ + x + x = ẋ + x + x ,
2 2 4 2 2 4 t=0
1 2 1 2 ǫ 4 1
ẋ + x + x = (2 + ǫ)
2 2 4 4
indicating that the solution is bounded. However, the series has a secular term −ǫ 83 t sin t that grows
without bound. This solution is only valid for t ≪ ǫ−1 .
Error
6 Numerical - O(1)
4
t
20 40 60 80 100
-2
-4
ε = 0.2
-6
Figure 4.12: Difference between exact and leading order solution to Duffing’s equation
The difference between the exact solution and the leading order solution, xexact (t) − (xo (t) + ǫx1 (t))
is plotted in Figure 4.13. There is some improvement for early time, but the solution is actually worse
for later time. This is because of the secularity.
To have a solution valid for all time, we strain the time coordinate
t = (1 + c1 ǫ + c2 ǫ2 + · · ·)τ
where τ is the new time variable. The ci ’s should be chosen to avoid secular terms.
Differentiating
−1
dx dτ dx dt
ẋ = =
dτ dt dτ dτ
dx
= (1 + c1 ǫ + c2 ǫ2 + · · ·)−1
dτ
d2 x
ẍ = (1 + c1 ǫ + c2 ǫ2 + · · ·)−2
dτ 2
d2 x
= (1 − 2((c1 ǫ + c2 ǫ2 + · · ·) + 3(c1 ǫ + c2 ǫ2 + · · ·)2 + · · ·))
dτ 2
d2 x
= (1 − 2c1 ǫ + (3c21 − 2c2 )ǫ2 + · · ·)
dτ 2
Furthermore, we write
x = x0 + ǫx1 + ǫ2 x2 + . . .
Error
t
20 40 60 80 100
-2
-4 ε = 0.2
-6
Figure 4.13: Difference between exact and uncorrected solution to O(ǫ) for Duffing’s equation
Collecting terms
d2 x0 dx0
O(ǫ0 ) : dτ 2 + x0 = 0, x0 (0) = 1, dτ (0) =0
x0 (τ ) = cos τ
d2 x1 2
O(ǫ1 ) : dτ 2 + x1 = 2c1 ddτx20 − x30 , x1 (0) = 0, dx1
dτ (0) =0
= −2c1 cos τ − cos3 τ
= −(2c1 + 34 ) cos τ − 41 cos 3τ
1
x1 (τ ) = 32 (− cos τ + cos 3τ ) if we choose c1 = − 83
Thus
1
x(τ ) = cos τ + ǫ (− cos τ + cos 3τ ) + · · ·
32
Since
3
t = 1 − ǫ + ··· τ
8
3
τ = 1 + ǫ + ··· t
8
so that
3 1 3 3
x(t) = cos 1 + ǫ + ··· t + ǫ − cos 1 + ǫ + · · · t + cos 3 1 + ǫ + · · · t +···
8 32 8 8
| {z }
Frequency Modulation (FM)
The difference between the exact solution and the leading order solution, xexact (t) − (xo (t) + ǫx1 (t))
for the corrected solution to O(ǫ) is plotted in Figure 4.14. The error is much smaller relative to the
previous cases; there does appear to be a slight growth in the amplitude of the error with time. This
might not be expected, but in fact is a characteristic behavior of the truncation error of the numerical
method used to generate the exact solution.
Error
6
Numerical - [O(1) + O(ε)]
4 Corrected
t
20 40 60 80 100
-2
ε = 0.2
-4
-6
Figure 4.14: Difference between exact and corrected solution to O(ǫ) for Duffing’s equation
Example 4.12
Find the amplitude of the limit cycle oscillations of the van der Pol equation
d2 x dx
(1 − 2c1 ǫ + . . .) − ǫ(1 − x2 ) (1 − c1 ǫ + . . .) + x = 0.
dτ 2 dτ
We also use
x = x0 + ǫx1 + ǫ2 x2 + . . . .
Thus we get
x0 = A cos τ
0
to O(ǫ ). To O(ǫ), the equation is
d2 x1 A2 A3
2
+ x1 = −2c 1 A cos τ − A 1 − sin τ + sin 3τ
dτ 4 4
x Numerical Solution
2
Method of
Strained Coordinates
1
2
x’’ - ε (1 - x ) + x = 0
t x(0) = 2, x’(0) = 0
10 20 30 40 50
-1 ε = 0.1
-2
Error Error
0.2 Numerical - O(1) 0.2
Numerical - [O(1) + O(ε)]
0.1 0.1
t t
10 20 30 40 50 10 20 30 40 50
-0.1 -0.1
-0.2 -0.2
Figure 4.15: Exact, difference between exact and asymptotic leading order solution, and
difference between exact and corrected asymptotic solution to O(ǫ) for van der Pol equation
for the method of strained coordinates
dx
dt
2
x
-2 -1 1 2
-1
-2
Figure 4.16: Phase plane solution trajectories for solution of van der Pol equation; d2 x/dt2 −
ǫ(1 − x2 )dx/dt + x = 0, x(0) = 2, dx/dt(0) = 0; ǫ = 1/10.
The exact solution xexact (t), the difference between the exact solution and the asymptotic leading
order solution, xexact (t) − xo (t), and the difference between the exact solution and the asymptotic
solution corrected to O(ǫ): xexact (t) − (xo (t) + ǫx1 (t)) is plotted in Figure 4.15.
Example 4.13
Solve
d2 x dx
2
− ǫ(1 − x2 ) + x = 0, with x(0) = 0, dx
dt (0) =1
dt dt
Let x = x(τ, τ̃ ), where the fast time scale is
τ = (1 + a1 ǫ + a2 ǫ2 + · · ·)t
d2 x 2
2
2∂ x 2 ∂2x 2
2∂ x
= (1 + a 1 ǫ + a 2 ǫ + · · ·) + 2(1 + a 1 ǫ + a 2 ǫ + · · ·)ǫ + ǫ
dt2 ∂τ 2 ∂τ ∂ τ̃ ∂ τ̃ 2
Introduce
x = x0 + ǫx1 + ǫ2 x2 + · · ·
So to O(ǫ), the differential equation becomes
∂ 2 x0 ∂x0
+ x0 = 0 with x0 (0, 0) = 0, ∂τ (0, 0) =1
∂τ 2
The solution is
x0 = A(τ̃ ) cos τ + B(τ̃ ) sin τ with A(0) = 0, B(0) = 1
The terms of O(ǫ1 ) give
∂ 2 x1 ∂ 2 x0 ∂ 2 x0 ∂x0
+ x1 = −2a1 − 2 + (1 − x20 )
∂τ 2 ∂τ 2 ∂τ ∂ τ̃ ∂τ
′ A 2 2
= 2a1 B + 2A − A + (A + B ) sin τ
4
B
+ 2a1 A − 2B ′ + B − (A2 + B 2 ) cos τ
4
A 2 B
+ (A − 3B 2 ) sin 3τ − (3A2 − B 2 ) cos 3τ
4 4
with
x1 (0, 0) = 0
∂x1 ∂x0 ∂x0
(0, 0) = −a1 (0, 0) − (0, 0)
∂τ ∂τ ∂ τ̃
∂x0
= −a1 − (0, 0)
∂ τ̃
Since ǫt is already represented in τ̃ , choose a1 = 0. Then
A 2
2A′ − A + (A + B 2 ) = 0
4
B
2B ′ − B + (A2 + B 2 ) = 0
4
Since A(0) = 0, try A(τ̃ ) = 0. Then
B3
2B ′ − B + =0
4
Multiplying by B, we get
B4
2BB ′ − B 2 + =0
4
B4
(B 2 )′ − B 2 + =0
4
Taking F ≡ B 2 , we get
F2
F′ − F + =0
4
This is a first order ODE in F , which can be easily solved. Separating variables, integrating, and
transforming from F back to B, we get
B2 τ̃
2 = Ce
1 − B4
2
B= √
1 + 3e−τ̃
x Numerical Solution
2
Method of Multiple Scales
1 x’’ - ε (1 - x2 ) x’ + x = 0
x(0) = 0, x’(0) = 1
t
10 20 30 40 50
ε = 0.1
-1
-2
Error Error
1 Numerical - O(1) 0.1
Numerical - [O(1) + O(ε)]
Corrected
0.5 0.05
t t
10 20 30 40 50 10 20 30 40 50
-0.5 -0.05
-1 -0.1
Figure 4.17: Exact, difference between exact and asymptotic leading order solution, and
difference between exact and corrected asymptotic solution to O(ǫ) for van der Pol equation
for the method of multiple scales
so that
2
x(τ, τ̃ ) = √ sin τ + O(ǫ)
1 + 3e−τ̃
2
x(t) = √ sin (1 + O(ǫ2 ))t + O(ǫ)
1 + 3e−ǫt
| {z }
Amplitude Modulation (AM)
The exact solution xexact (t), the difference between the exact solution and the asymptotic leading
order solution, xexact (t) − xo (t), and the difference between the exact solution and the asymptotic
solution corrected to O(ǫ): xexact (t) − (xo (t) + ǫx1 (t)) is plotted in Figure 4.17.
Note that the amplitude, which is initially 1, grows to a value of 2, the same value which was
obtained in the previous example. Here, we have additionally obtained the time scale for the growth
of the amplitude change. Note also that the leading order approximation is quite poor for t > 1ǫ , while
the corrected approximation is relatively quite good.
dx
dt
2
x
-2 -1 1 2
-1
-2
Figure 4.18: Phase plane solution trajectories for solution of van der Pol equation; d2 x/dt2 −
ǫ(1 − x2 )dx/dt + x = 0, x(0) = 0, dx/dt(0) = 1; ǫ = 1/10.
Example 4.14
Solve
ǫy ′′ + y ′ + y = 0, with y(0) = 0, y(1) = 1 (4.11)
An exact solution to this equation exists, namely
√
sinh x 1−4ǫ
1−x 2ǫ
y(x) = exp √
2ǫ sinh 1−4ǫ
2ǫ
We could in principle simply expand this in a Taylor series in ǫ. However, for more difficult problems,
exact solutions are not available. So here we will just use the exact solution to verify the validity of the
method.
We begin with a regular perturbation expansion
The solution obtained above is the solution in the outer region. To satisfy the boundary condition
y0 (1) = 1, we find that a = e, so that
y = e1−x + · · ·
In the inner region, we choose a new independent variable X defined as X = x/ǫ, so that the equation
becomes
d2 y dy
2
+ + ǫy = 0
dX dX
Using a perturbation expansion, the lowest order equation is
d2 y0 dy0
2
+ =0
dX dX
with a solution
y0 = A + Be−X
Applying the boundary condition y0 (0) = 0, we get
y0 = A(1 − e−X )
yinner (X → ∞) = youter (x → 0)
lim y = e
x→∞
and
y(x) = e1−x + · · · in the outer region
lim y = e
x→0
A composite solution can also be written by adding the two solutions and subtracting the common
part.
y(x) = e(1 − e−x/ǫ ) + · · · + e1−x + · · · − e
y = e(e−x − e−x/ǫ ) + · · ·
The exact solution, the inner layer solution, the outer layer solution, and the composite solution
are plotted in Figure 4.19.
Example 4.15
Obtain the solution of the previous problem
Outer Layer
y Solution
ε y’’ + y’ + y = 0
2.5
Inner Layer y (0) = 0
Exact Solution
2 Solution y (1) = 1
ε = 0.1
1.5
1 Composite Prandtl’s
Solution Boundary Layer Method
0.5
x
0.2 0.4 0.6 0.8 1
Figure 4.19: Exact, inner layer solution, outer layer solution, and composite solution for
boundary layer problem
Figure 4.20: Difference between exact and asymptotic solutions for two different orders of
approximation for a boundary layer problem
and
y(x) = e1−x + ǫ(1 − x)e1−x · · · in the outer region
The composite solution is
y = e1−x − (1 + x)e1−x/ǫ + ǫ (1 − x)e1−x − e1−x/ǫ + · · ·
The difference between the exact solution and the approximation from the previous example, and the
difference between the exact solution and approximation from this example are plotted in Figure 4.20.
Example 4.16
In the same problem, investigate the possibility of having the boundary layer at x = 1. The outer
solution now satisfies the condition y(0) = 0, giving y = 0. Let
x−1
X=
ǫ
The lowest order inner solution satisfying y(X = 0) = 1 is
y = A + (1 − A)e−X
However, as X → −∞, this becomes unbounded and cannot be matched with the outer solution. Thus,
a boundary layer at x = 1 is not possible.
y Error
1 ε y’’ - y’ + y = 0 0.1
0.2 0.02
x x
0 0.2 0.4 0.6 0.8 1 0 0.2 0.4 0.6 0.8 1
Approximate
Figure 4.21: Exact, approximate, and difference in predictions for a boundary layer problem
Example 4.17
Solve
ǫy ′′ − y ′ + y = 0, with y(0) = 0, y(1) = 1
The boundary layer is at x = 1. The outer solution is y = 0. Taking
x−1
X=
ǫ
the inner solution is
y = A + (1 − A)eX + . . .
Matching, we get
A=0
so that we have a composite solution
y(x) = e(x−1)/ǫ + . . .
The exact solution, the approximate solution to O(ǫ), and the difference between the exact solution
and the approximation, are plotted in Figure 4.21.
where
Z
1 x
v(x) = y(x) exp − P (s)ds (4.15)
2 0
1 dP 1
R(x) = Q(x) − − (P (x))2 (4.16)
2 dx 4
So it is sufficient to study equations of the form (4.14). The Wentzel,7 Kramers,8 Brillouin,9
(WKB) method is used for equations of the kind
d2 y
ǫ2 = f (x)y (4.17)
dx2
where ǫ is a small parameter. This also includes an equation of the type
d2 y
ǫ2 = (λ2 p(x) + q(x))y (4.18)
dx2
where λ is a large parameter. Alternatively, by taking x = ǫt, equation (4.17) becomes
d2 y
= f (ǫt)y (4.19)
dt2
We can also write equation (4.17) as
d2 y
= g(x)y (4.20)
dx2
where g(x) is slowly varying in the sense that g ′/g 3/2 ∼ O(ǫ).
We seek solutions to equation (4.17) of the form
Z x
1 2
y(x) = exp (S0 (s) + ǫS1 (s) + ǫ S2 (s) + · · ·)ds (4.21)
ǫ xo
from which p
S0 (x) = ± f (x)
To O(ǫ1 ) we have
dSo
2S0 (x)S1 (x) + =0
dx
from which
dSo
dx
S1 (x) = −
2So (x)
df
± √1 dx
2 f (x)
S1 (x) = − p
2 ± f (x)
df
dx
S1 (x) = −
4f (x)
Thus we get the general solution
Z x
1
y(x) = C1 exp (S0 (s) + ǫS1 (s) + · · ·)ds (4.22)
ǫ xo
Z x
1
+C2 exp (S0 (s) + ǫS1 (s) + · · ·)ds (4.23)
ǫ xo
Z !
df
1 x p
y(x) = C1 exp ( f (s) − ǫ ds + · · ·)ds (4.24)
ǫ xo 4f (s)
Z !
df
1 x p
+C2 exp (− f (s) − ǫ ds + · · ·)ds (4.25)
ǫ xo 4f (s)
Z f (x) ! Z x
df 1 p
y(x) = C1 exp − exp ( f (s) + · · ·)ds (4.26)
f (xo ) 4f ǫ xo
Z f (x) ! Z
df 1 x p
+C2 exp − exp − ( f (s) + · · ·)ds (4.27)
f (xo ) 4f ǫ xo
Z x Z
Ĉ1 1 p Ĉ2 1 xp
y(x) = exp f (s)ds + exp − f (s)ds + · (4.28)
··
(f (x))1/4 ǫ xo (f (x))1/4 ǫ xo
This solution is not valid near x = a for which f (a) = 0. These are called turning points.
Example 4.18
Find an approximate solution of the Airy10 equation
ǫ2 y ′′ + xy = 0, for x > 0
In this case
f (x) = −x
so that √
S0 (x) = ±i x
and
S0′ 1
S1 (x) = − =−
2S0 4x
The solutions are of the form
Z Z
i √ dx
y = exp ± x dx − + ···
ǫ 4x
1 2x3/2 i
= exp ± + ···
x1/4 3ǫ
Here Ai and Bi are Airy functions of the first and second kind, respectively.
Example 4.19
Find a solution of x3 y ′′ = y, for small, positive x.
Let ǫ2 X = x, so that X is of O(1) when x is small. Then the equation becomes
d2 y
ǫ2 = X −3 y
dX 2
The WKB method is applicable. We have f = X −3 . The general solution is
′ 3/4 2 ′ 3/4 2
y = C1 X exp − √ + C2 X exp √ + ···
ǫ X ǫ X
10
George Biddell Airy, 1801-1892, English applied mathematician, First Wrangler at Cambridge, holder
of the Lucasian Chair (that held by Newton) at Cambridge, Astronomer Royal who had some role in delaying
the identification of Neptune as predicted by John Couch Adams’ perturbation theory in 1845.
Here I1 is a modified Bessel function of the first kind of order one, and K1 is a modified Bessel function
of the second kind of order one.
Example 4.20
Solve
x3 y ′′ = y
for small, positive x. Let y = eS(x) , so that y ′ = S ′ eS , y ′′ = (S ′ )2 eS + S ′′ eS , from which
S ′′ + (S ′ )2 = x−3
Assume that S ′′ ≪ (S ′ )2 (to be checked later). Thus S ′ = ±x−3/2 , and S = ±2x−1/2 . Checking we
get S ′′ /(S ′ )2 = x1/2 → 0 as x → 0, confirming the assumption. Now we add a correction term so that
S(x) = 2x−1/2 + C(x), where we have taken the positive sign. Assume that C ≪ 2x−1/2 . Substituting
in the equation, we have
3 −5/2
x + C ′′ − 2x−3/2 C ′ + (C ′ )2 = 0
2
Since C ≪ 2x−1/2 , we have C ′ ≪ x−3/2 and C ′′ ≪ 32 x−5/2 . Thus
3 −5/2
x − 2x−3/2 C ′ = 0
2
from which C ′ = 34 x−1 and C = 43 ln x. We can now check the assumption on C.
We have S(x) = 2x−1/2 + 43 ln x, so that
2
y = x3/4 exp − √ + ···
x
Example 4.21
Solve
y ′ = e−xy
for y > 0 and large x.
As x → ∞, y ′ → 0, so that y → c (a positive constant). Substituting into the equation, we have
y ′ = e−cx
Problems
1. Solve as a series in x for x > 0 about the point x = 0:
In each case find the exact solution with a symbolic computation program, and compare graphically
the first four terms of your series solution with the exact solution.
2. Find two term expansions for each of the roots of
where λ is large.
First
y Approximation, y = 1 - exp(-x)
1
0.8
8
0.2
Repeated Substitution Method
x
2 4 6 8 10
-0.2
-0.4
Figure 4.22: Numerical and first approximate solution for repeated substitution problem
λ
y ′′ + y=0
λ+x
with y(0) = 0, y(1) = 1, where λ is a large parameter. For λ = 20, plot y(x) for the two term
expansion. Also compute the exact solution by numerical integration. Plot the difference between the
asymptotic and numerical solution versus x.
4. Find the leading order solution for
dy
(x − ǫy) + xy = e−x ,
dx
where y(1) = 1, and x ∈ [0, 1], ǫ ≪ 1. For ǫ = 0.2, plot the asymptotic solution, the exact solution
and the difference versus x.
5. The motion of a pendulum is governed by the equation
d2 x
+ sin(x) = 0
dt2
with x(0) = ǫ, dx
dt (0) = 0. Using strained coordinates, find the approximate solution of x(t) for small ǫ
through O(ǫ2 ). Plot your results for both your asymptotic results and those obtained by a numerical
integration of the full equation.
6. Find an approximate solution for
y ′′ − yey/10 = 0,
with y(0) = 1, y(1) = e.
11. Find all solutions through O(ǫ2 ), where ǫ is a small parameter, and compare with the exact result for
ǫ = 0.01.
where ǫ is a small parameter. Compare graphically with the exact solution for ǫ = 0.3 and 0 ≤ t ≤ 2.
14. Write down an approximation for
Z π/2 p
1 + ǫ cos2 x dx
0
for small ǫ. Compare the approximate and exact solutions graphically in the range 0 ≤ x ≤ 1 for (a)
ǫ = 0.1, (b) ǫ = 0.25, and (c) ǫ = 0.5.
19. Find an approximate solution to
for small, positive ǫ. Compare with the exact solution. Plot both the exact solution and the approxi-
mate solution on the same graph for A = 1, B = 0, ǫ = 0.3.
20. Find an approximate solution to the following problem for small ǫ
where λ ≫ 1.
24. Find the complementary functions of
y ′′′ − xy = 0
in terms of expansions near x = 0. Retain only two terms for each function.
25. Find, correct to O(ǫ), the solution of
y ′′ + ǫy 2 = 0
with y(0) = 0, y(1) = 1 for ǫ ≪ 1. For ǫ = 2.5, compare the O(1), O(ǫ), and O(ǫ2 ) solutions to a
numerically obtained solution in x ∈ [0, 1].
28. Obtain a power series solution (in summation form) for y ′ + ky = 0 about x = 0, where k is an
arbitrary, nonzero constant. Compare to a Taylor series expansion of the exact solution.
29. Obtain two terms of an approximate solution for ǫex = cos x when ǫ is small. Graphically compare
to the actual values (obtained numerically) when ǫ = 0.2, 0.1, 0.01.
30. Obtain three terms of a perturbation solution for the roots of the equation (1 − ǫ)x2 − 2x + 1 = 0.
(Hint: The expansion x = x0 + ǫx1 + ǫ2 x2 + . . . will not work.)
31. The solution of the matrix equation A · x = y can be written as x = A−1 · y. Find the nth term of
the perturbation solution of (A + ǫB) · x = y, where ǫ is a small parameter. Obtain the first three
terms of the solution for
1 2 1 1/10 1/2 1/10 1/2
A = 2 2 1, B = 0 1/5 0 , y = 1/5 .
1 2 3 1/2 1/10 1/2 1/10
32. Obtain leading and first order terms for u and v, governed by the following set of coupled differential
equations, for small ǫ:
d2 u du 1 1
2
+ ǫv = 1, u(0) = 0, u(1) = + ǫ
dx dx 2 120
d2 v dv 1 1
2
+ ǫu = x, v(0) = 0, v(1) = + ǫ
dx dx 6 80
Compare asymptotic and numerically obtained results for ǫ = 0.2.
33. Obtain two terms of a perturbation solution to ǫfxx + fx = −e−x with boundary conditions f (0) = 0,
f (1) = 1. Graph the solution for ǫ = 0.2, 0.1, 0.05, 0.025 on 0 ≤ x ≤ 1.
34. Find two uniformly valid approximate solutions of
ω2u
ü + = 0 with u(0) = 0
1 + u2
up to the first order. Note that ω is not small.
35. Using a two-variable expansion, find the lowest order solution of
(a) ẍ + ǫẋ + x = 0 with x(0) = 0, ẋ(0) = 1
(b) ẍ + ǫẋ3 + x = 0 with x(0) = 0, ẋ(0) = 1
where ǫ ≪ 1. Compare asymptotic and numerically obtained results for ǫ = 0.01.
36. Obtain a three-term solution of
where ǫ ≪ 1.
37. Find an approximate solution to the following problem for small ǫ
valid near x = 0.
46. A bead can slide along a circular hoop in√a vertical plane. The bead is initially at the lowest position,
θ = 0, and given an initial velocity of 2 gR, where g is the acceleration due to gravity and R is the
radius of the hoop. If the friction coefficient is µ, find the maximum angle θmax reached by the bead.
Compare perturbation and numerical results. Present results on a θmax vs. µ plot, for 0 ≤ µ ≤ 0.3.
47. The initial velocity downwards of a body of mass m immersed in a very viscous fluid is V . Find
the velocity of the body as a function of time. Assume that the viscous force is proportional to the
velocity. Assume that the inertia of the body is small, but not negligible, relative to viscous and
gravity forces. Compare perturbation and exact solutions graphically.
48. For small ǫ, solve to lowest order using the method of multiple scales
λ
y ′′ + y=0
λ+x
with y(0) = 0, y(1) = 1, where λ is a large parameter.
53. Find all solutions of eǫx = x2 through O(ǫ2 ), where ǫ is a small parameter.
54. Solve
(1 + ǫ)y ′′ + ǫy 2 = 1
with y(0) = 0, y(1) = 1 through O(ǫ2 ), where ǫ is a small parameter.
55. Solve to lowest order
ǫy ′′ + y ′ + ǫy 2 = 1
with y(0) = −1, y(1) = 1, where ǫ is a small parameter. For ǫ = 0.2, plot asymptotic and numerical
solutions to the full equation.
56. Find the series solution of the differential equation
y ′′ + xy = 0
sin x = ǫ cos 2x
(1 + ǫy)y ′′ + ǫy ′2 − N 2 y = 0
with y ′ (0) = 0, y(1) = 1. for small ǫ. N is a constant. Plot the asymptotic and numerical solution for
ǫ = 0.12, N = 10.
61. Solve
1
ǫy ′′ + y ′ =
2
with y(0) = 0, y(1) = 1 for small ǫ. Plot asymptotic and numerical solutions for ǫ = 0.12.
62. Find if the van der Pol equation
ÿ − ǫ(1 − y 2 )ẏ + k 2 y = 0
has a limit cycle of the form y = A cos ωt.
63. Solve y ′ = e−2xy for large x where y is positive. Plot y(x).
Solutions of differential equations give rise to complementary functions. Some of these are
well known such as sin and cos. This chapter will consider these and other functions which
arise from the solution from a variety of second order differential equations with non-constant
coefficients. The notion of eigenvalues, eigenfunctions, and orthogonal functions will be
introduced; a stronger foundation will be built in the chapter on linear analysis. It will be
shown how one can expand an arbitrary function in terms of infinite sums of the product of
scalar amplitudes with orthogonal basis functions.
137
138 CHAPTER 5. ORTHOGONAL FUNCTIONS AND FOURIER SERIES
Z x
c(x) b(s)
q(x) = exp ds (5.6)
a(x) xo a(s)
With these definitions, the original equation is transformed to the type known as a
Sturm-Liouville1 equation:
d dy
p(x) + (q(x) + λr(x)) y(x) = 0 (5.7)
dx dx
1 d d
p(x) + q(x) y(x) = −λ y(x) (5.8)
r(x) dx dx
Here the Sturm-Liouville linear operator Ls is
1 d d
Ls = p(x) + q(x)
r(x) dx dx
so we have
Ls y(x) = −λ y(x)
Now the trivial solution y(x) = 0 will satisfy the differential equation. In addition for special
values of λ, known as eigenvalues, there are certain functions, known as eigenfunctions which
also satisfy the differential equation.
Now it can be shown that if we have for x ∈ [x0 , x1 ]
then an infinite number of real positive eigenvalues λ and corresponding eigenfunctions yi (x)
exist for which the differential equation is satisfied. Moreover it can also be shown (Hilde-
brand, p. 204) that a consequence of the homogeneous boundary conditions is the orthogo-
nality condition:
Z x1
r(x)yi (x)yj (x) dx = 0, for i 6= j (5.12)
x0
Z x1
r(x)yi (x)yi (x) dx = K 2 (5.13)
x0
Here K is a real constant. This can be written compactly using the Kronecker delta function,
δij as
Z x1
r(x)yi (x)yj (x) dx = K 2 δij . (5.14)
x0
1
Jacques Charles François Sturm, 1803-1855, Swiss-born French mathematician and Joseph Liouville,
1809-1882, French mathematician.
Sturm-Liouville theory shares many analogies with vector algebra. In the same sense that
the dot product of a vector with itself is guaranteed positive, we wish to define a “product”
for the eigenfunctions in which the “product” of a function and itself is guaranteed positive.
Consequently, the eigenfunctions of a Sturm-Liouville operator Ls are said to be orthogonal
to each other.
Based on the above result we can define functions ϕi (x):
r
r(x)
ϕi (x) = yi (x) (5.15)
K2
so that Z x1
ϕi (x)ϕj (x) dx = δij . (5.16)
x0
Such functions are said to be orthonormal. While orthonormal functions have great
utility, note that in the context of our Sturm-Liouville nomenclature, that ϕi (x) does not in
general satisfy the Sturm-Liouville equation: Ls ϕi (x) 6= −λi ϕi (x). If, however, r(x) = C,
where C is a scalar constant, then in fact Ls ϕi (x) = −λi ϕi (x). Whatever the case, in all
cases we are guaranteed Ls yi (x) = −λi yi (x). The yi (x) functions are orthogonal under the
influence of the weighting function r(x), but not necessarily orthonormal.
d2 y
+ λy = 0 (5.17)
dx2
dy
αy(a) + β (a) = 0 (5.18)
dx
dy
γy(b) + δ (b) = 0 (5.19)
dx
Here we have
a(x) = 1 (5.20)
b(x) = 0 (5.21)
c(x) = 0 (5.22)
so
p(x) = 1 (5.23)
r(x) = 1 (5.24)
q(x) = 0 (5.25)
So we can consider the domain x ∈ (−∞, ∞). In practice it is more common to consider the
finite domain in which x ∈ [x0 , x1 ]. The Sturm-Liouville operator is
d2
Ls =
dx2
The eigenvalue problem is
d2
y(x) = −λ y(x)
dx2
The general solution is √ √
y(x) = C1 cos( λx) + C2 sin( λx)
Example 5.1
Find the eigenvalues and eigenfunctions for
d2 y
+ λy = 0
dx2
y(0) = y(ℓ) = 0
For y(0) = 0 we get
√ √
y(0) = 0 = C1 cos( λ(0)) + C2 sin( λ(0))
C1 = 0
So √
y(x) = C2 sin( λx)
At the other boundary we have √
y(ℓ) = 0 = C2 sin( λ(ℓ))
For non-trivial solutions we need C2 6= 0, which then requires that
√
λℓ = nπ n = ±1, ±2, ±3, . . .
so nπ 2
λ=
ℓ
The eigenvalues and eigenfunctions are
n2 π 2
λn =
ℓ2
and nπx
yn (x) = sin
ℓ
respectively.
Legendre’s2 equation is below. Here, it is convenient to let the term n(n + 1) play the
role of λ.
d2 y dy
(1 − x2 ) 2 − 2x + n(n + 1)y = 0 (5.34)
dx dx
Here
a(x) = 1 − x2 (5.35)
b(x) = −2x (5.36)
c(x) = 0 (5.37)
2
Adrien-Marie Legendre, 1752-1833, French/Parisian mathematician.
r(x) = 1 (5.42)
q(x) = 0 (5.43)
Pn (x) P4 P3
3 P
2
2
P
1
1 P
0
-2
Direct substitution shows that Pn (x) satisfies both the differential equation and the orthog-
onality condition. It is then easily shown that the following functions are orthonormal on
the interval x ∈ (−1, 1):
r
1
ϕn (x) = n + Pn (x) (5.59)
2
T (x)
n
2 T1
1.5
1 T0
0.5
-2 -1 1 2
x
-0.5
-1
-1.5
T3 T
T 4
2
-2
Thus
√ d √ d
Ls = 1 − x2 1 − x2 (5.73)
dx dx
That the two forms are equivalent can be easily checked by direct expansion of the above
equation.
The first five eigenfunctions of the Chebyshev equation are plotted in the Figure 5.2.
They can be expressed in terms of polynomials known as the Chebyshev polynomials, Tn (x).
These polynomials can be obtained by a regular series expansion of the original differential
equation.
Eigenvalues and eigenfunctions are listed below:
The Rodrigues4 formula gives a generating formula for general n. The orthogonality condi-
4
Benjamin Olinde Rodrigues, 1794-1851, obscure French mathematician, of Portuguese and perhaps
Spanish roots.
tion is
Z 1
Ti (x)Tj (x)
√ dx = 0 i 6= j (5.81)
−1 1 − x2
Z 1
Ti (x)Ti (x) π if i = j = 0
√ dx = π (5.82)
−1 1 − x2 2
if i = j = 1, 2, . . .
Direct substitution shows that Tn (x) satisfies both the differential equation and the orthog-
onality condition. We can deduce then that the functions ϕn (x)
q
√1 Tn (x) if n = 0
π 1−x2
ϕn (x) = q (5.83)
√2 Tn (x) if n = 1, 2, . . .
π 1−x 2
or Z 1
ϕi (x)ϕj (x) = δij . (5.85)
−1
H4 Hn(x)
30 H2
20
10
H0 H1
-3 -2 -1 1 2 3 x
-10
-20
-30
H3
So
d
x2 −x2 d
Ls = e e (5.92)
dx dx
The first five eigenfunctions of the Hermite equation are plotted in the Figure 5.3. They
can be expressed in terms of polynomials known as the Hermite polynomials, Hn (x). These
polynomials can be obtained by a regular series expansion of the original differential equation.
Eigenvalues and eigenfunctions are listed below:
Direct substitution shows that Hn (x) satisfies both the differential equation and the orthog-
onality condition. It is then easily shown that the following functions are orthonormal on
the interval (−∞, ∞)
2
e−x /2 Hn (x)
ϕn (x) = p√ (5.102)
π2n n!
So
d
x −x d
Ls = e xe (5.109)
dx dx
The first five eigenfunctions of the Laguerre equation are plotted in the Figure 5.4. They
can be expressed in terms of polynomials known as the Laguerre polynomials, Ln (x). These
polynomials can be obtained by a regular series expansion of the original differential equation.
Eigenvalues and eigenfunctions are listed below:
L n (x) L2
L4
10
5
L0
-2 2
x
4 6 8 10
-5
-10 L1
L 3
3 1
λ=3 L3 (x) = 1 − 3x + x2 − x3 (5.113)
2 6
2 3 1
λ=4 2
L4 (x) = 1 − 4x + 3x − x + x4 (5.114)
3 24
..
. (5.115)
n n −x
1 x d (x e )
λ=n Ln (x) = e Rodrigues’ formula (5.116)
n! dxn
The orthogonality condition reduces to
Z ∞
e−x Li (x)Lj (x) dx = δij (5.117)
0
Direct substitution shows that Ln (x) satisfies both the differential equation and the orthog-
onality condition. It is then easily shown that the following functions are orthonormal on
the interval x ∈ (0, ∞):
where Jν (µx) and Yν (µx) are called the Bessel and Neumann functions of order ν. Often
Jν (µx) is known as a Bessel function of the first kind and Yν (µx) is known as a Bessel
function of the second kind. Both Jν and Yν are represented by infinite series rather than
finite series such as the series for Legendre polynomials.
The Bessel function of the first kind of order ν, Jν (µx), is represented by
ν X∞ k
1 − 14 µ2 x2
Jν (µx) = µx (5.133)
2 k=0
k!Γ(ν + k + 1)
The Neumann function Yν (µx) has a complicated series representation (see Hildebrand).
The representations for J0 (µx) and Y0 (µx) are
1 2 2 1
1 2 2 2
1 2 2 n
µ x µ x − µ x
J0 (µx) = 1 − 4 + 4 + ...+ 4
(1!)2 (2!)2 (n!)2
2 1
Y0 (µx) = ln µx + γ J0 (µx)
π 2
1 2 2 1
1 2 2 2 !
2 4
µ x 1 4
µx
+ 2
− 1+ ...
π (1!) 2 (2!)2
The Bessel functions J0 (x), J1 (x), J2 (x), J3 (x), and J4 (x) along with the Neumann functions
Y0 (x), Y1 (x), Y2 (x), Y3 (x), and Y4 (x) are plotted in Figure 5.6 (so here µ = 1).
The orthogonality condition for a domain x ∈ (0, 1), taken here for the case in which the
eigenvalue is µi , can be shown to be
Z 1
xJν (µi x)Jν (µj x) dx = 0 i 6= j (5.137)
0
Z 1
1
xJν (µi x)Jν (µi x) dx = (Jν+1 (µi ))2 i=j (5.138)
0 2
Here we must choose µi such that Jν (µi ) = 0, which corresponds to a vanishing of the
function at the outer limit x = 1. See Hildebrand, p. 226.
Jo(µnx)
1
Jo(µ0x)
0.8
0.6
0.4
0.2
-0.4
Jo(µ3x) Jo(µ2x) Jo(µ1x)
Figure 5.5: Bessel functions J0 (µ0 x), J0 (µ1 x), J0 (µ2 x), J0 (µ3 x)
Y (x)
J (x)
ν ν
1 J 1
0
0.75
0.8 Y0
J1 0.5 Y1 Y2
0.6
Y3 Y4
J2 J3 J4 0.25
0.4
x
2 4 6 8 10
0.2
-0.25
x
2 4 6 8 10 -0.5
-0.2 -0.75
-0.4 -1
Figure 5.6: Bessel functions J0 (x), J1 (x), J2 (x), J3 (x), J4 (x) and Neumann functions Y0 (x),
Y1 (x), Y2 (x), Y3 (x), Y4 (x)
d2 y dy
x2 2
+ x − (x2 + ν 2 )y = 0 (5.142)
dx dx
the solutions of which are the modified Bessel functions. It is satisfied by the modified Bessel
functions. The modified Bessel function of the first kind of order ν is
d2 y dy
x2 + x − (p2 + ix2 )y = 0 (5.145)
dx2 dx
where p is a real constant, are called the ber and bei functions.
We generally truncate the infinite series to a finite number of N terms so that f (x) is
approximated by
XN
f (x) ∼ an ϕn (x). (5.147)
n=1
From hereout, we will use an equality for the N-term approximation, while realizing that
we are actually tolerating an error. The represenation is useful only if the infinite series
converges so that the error incurred in neglecting terms past N are small relative to the
terms included. The problem is to determine what the coefficients an must be. They can
be found in the following manner. We first assume the expansion exists and multiply both
sides by ϕk (x):
N
X
f (x)ϕk (x) = an ϕn (x)ϕk (x) (5.148)
n=0
Z x1 Z N
x1 X
f (x)ϕk (x) dx = an ϕn (x)ϕk (x) dx (5.149)
x0 x0 n=0
N
X Z x1
= an ϕn (x)ϕk (x) dx (5.150)
n=0 x0
XN
= an δnk (5.151)
n=0
= a0 δ0k +a1 δ1k + . . . + ak δkk + . . . + aN δN k (5.152)
|{z} |{z} |{z} |{z}
=0 =0 =1 =0
= ak (5.153)
So trading k and n Z x1
an = f (x)ϕn (x) dx (5.154)
x0
The series is known as a Fourier8 series. Depending on the expansion functions, the series
is often specialized as Fourier-sine, Fourier-cosine, Fourier-Legendre, Fourier-Bessel, etc. We
have inverted Eq. (5.146) to solve for the unknown an . The inversion was aided greatly
by the fact that the basis functions were orthonormal. For non-orthonormal, as well as
non-orthogonal bases, more general techniques exist for the determination of an .
Example 5.2
Represent
f (x) = x2 on x ∈ [0, 3] (5.155)
with a series of
8
Jean Baptiste Joseph Fourier, 1768-1830, French mathematician.
• trigonometric functions
• Legendre polynomials
• Chebyshev polynomials
• Bessel functions
Trigonometric Series
For the trigonometric series let’s try a Fourier sine series. The orthonormal functions in this case
are r
2 nπx
ϕn (x) = sin (5.156)
3 3
The coefficients are thus
r Z 3 nπx
2
an = x2 sin dx (5.157)
3 0 3
so
a0 = 0 (5.158)
a1 = 4.17328 (5.159)
a2 = −3.50864 (5.160)
a3 = 2.23376 (5.161)
a4 = −1.75432 (5.162)
a5 = 1.3807 (5.163)
Note that the magnitude of the coefficent on the orthonormal function, an , decreases as n increases.
From this, one can loosely infer that the higher frequency modes contain less “energy.”
r πx
2 2πx
f (x) = 4.17328 sin − 3.50864 sin
3 3 3
3πx 4πx 5πx
+2.23376 sin − 1.75432 sin + 1.3807 sin + ...
3 3 3
The function f (x) = x2 and five terms are plotted in Figure 5.7
Legendre polynomials
Next let’s try the Legendre polynomials. The Legendre polynomials are orthogonal on x ∈ [−1, 1],
and we have x ∈ [0, 3], so let’s define
2
x̃ = x−1 (5.164)
3
3
x= (x̃ + 1) (5.165)
2
so that the domain x ∈ [0, 3] maps into x̃ ∈ [−1, 1]. So expanding x2 on the domain x ∈ [0, 3] is
equivalent to expanding
2
3 9
(x̃ + 1)2 = (x̃ + 1)2 x̃ ∈ [−1, 1] (5.166)
2 4
f(x)
x2
Fourier-sine series
8 (five terms)
x
0.5 1 1.5 2 2.5 3
Now r
1
ϕn (x̃) = n + Pn (x̃) (5.167)
2
So
r Z 1
9 1
an = n+ (x̃ + 1)2 Pn (x̃) dx̃ (5.168)
4 2 −1
Evaluating we get
√
a0 = 3 2 = 4.24264 (5.169)
r
3
a1 = 3 = 3.67423 (5.170)
2
3
a2 = √ = 0.948683 (5.171)
10
a3 = 0 (5.172)
..
. (5.173)
an = 0 n>3 (5.174)
Once again, the fact the a0 > a1 > a2 indicates the bulk of the “energy” is contained the lower frequency
modes. Carrying out the multiplication and returning to x space gives the finite series, which can be
expressed in a variety of forms:
2 !
9 2 3 1 3 2
= 3(1) + x−1 + − + x−1 , (5.178)
2 3 2 2 2 3
9 3
= 3 + − + 3x + − 3x + x2 , (5.179)
2 2
= x2 . (5.180)
Thus, the Fourier-Legendre representation is exact over the entire domain. This is because the function
which is being expanded has the same general functional form as the Legendre polynomials; both are
polynomials.
Chebyshev polynomials
Let’s now try the Chebyshev polynomials. These are orthogonal on the same domain as the Leg-
endre polynomials, so let’s use the same transformation as before.
Now
s
1
ϕn (x̃) = √ Tn (x̃) n=0 (5.181)
π 1 − x̃2
s
2
ϕn (x̃) = √ Tn (x̃) n>0 (5.182)
π 1 − x̃2
So
Z s
1
9 2 1
a0 = (x̃ + 1) √ T0 (x̃) dx̃ (5.183)
4 −1 π 1 − x̃2
Z s
1
9 2
an = (x̃ + 1)2 √ Tn (x̃) dx̃ (5.184)
4 −1 π 1 − x̃2
Evaluating we get
a0 = 4.2587 (5.185)
a1 = 3.4415 (5.186)
a2 = −0.28679 (5.187)
a3 = −1.1472 (5.188)
..
. (5.189)
With this representation, we see that |a3 | > |a2 |, so it is not yet clear that the “energy” is concentrated
in the high frequency modes. Consideration of more terms would verify that in fact it is the case that
the “energy ” of high frequency modes is decaying; in fact a4 = −0.683, a5 = −0.441, a6 = −0.328,
a7 = −0.254. So
v
u 2 4.2587 2 2
2 u
f (x) = x = t q √ T0 x − 1 + 3.4415 T1 x−1 (5.190)
2 2 3 3
π 1 − 2x − 1 3
2 2
−0.28679 T2 x − 1 − 1.1472 T3 x − 1 + ... (5.191)
3 3
The function f (x) = x2 and four terms are plotted in Figure 5.8
f(x)
10
8
Fourier-Chebyshev series x2
(four terms)
6
x
0.5 1 1.5 2 2.5 3
Bessel functions
Now let’s expand in terms of Bessel functions. The Bessel functions have been defined such that
they are orthogonal on a domain between zero and one when the eigenvalues are the zeros of the Bessel
function. To achieve this we adopt the transformation (and inverse):
x
x̃ = x = 3x̃.
3
With this transformation our domain transforms as follows:
x ∈ [0, 3] −→ x̃ ∈ [0, 1]
Now, the eigenvalues µn are such that J0 (µn ) = 0. We find using trial and error methods that solutions
for all the zeros can be found:
µ0 = 2.40483 (5.192)
µ1 = 5.52008 (5.193)
µ2 = 8.65373 (5.194)
..
. (5.195)
f (x)
x2
Fourier-Bessel Series
8
(ten terms)
6
Similar to the other functions, we could expand in terms of the orthonormalized Bessel functions, ϕn (x).
Instead, for variety, let’s directly operate on the above expression to determine the values for an .
∞
X
9x̃2 x̃J0 (µk x̃) = an x̃J0 (µn x̃)J0 (µk x̃) (5.196)
n=0
Z 1 Z 1X ∞
9x̃3 J0 (µk x̃) dx̃ = an x̃J0 (µn x̃)J0 (µk x̃) dx̃ (5.197)
0 0 n=0
Z 1 ∞
X Z 1
9 x̃3 J0 (µk x̃) dx̃ = an x̃J0 (µn x̃)J0 (µk x̃) dx̃ (5.198)
0 n=0 0
Z 1
= ak x̃J0 (µk x̃)J0 (µk x̃) dx̃ (5.199)
0
a0 = 4.44557 (5.201)
a1 = −8.3252 (5.202)
a2 = 7.2533 (5.203)
..
. (5.204)
Because the basis functions are not normalized, it is difficult to infer how the amplitude is decaying by
looking at an alone.
The function f (x) = x2 and ten terms of the Fourier-Bessel series approximation are plotted in
Figure 5.9 The Fourier-Bessel approximation is
x
f (x) = x2 = 4.44557 J0 2.40483 (5.205)
3x
−8.3252 J0 5.52008 (5.206)
x 3
+7.2533 J0 8.65373 + ... (5.207)
3
Note that other Fourier-Bessel expansions exist. Also note that even though the Bessel function does
not match the function itself at either boundary point, that the series still appears to be converging.
Problems
1. Show that oscillatory solutions of the delay equation
dx
(t) + x(t) + bx(t − 1) = 0
dt
are possible only when b = 2.2617. Find the frequency.
2. Show that xa Jν (bxc ) is a solution of
2a − 1 ′ a 2 − ν 2 c2
y ′′ − y + b2 c2 x2c−2 + y=0
x x2
Hence solve in terms of Bessel functions:
d2 y
(a) dx2 + k 2 xy = 0
d2 y
(b) dx2 + x4 y = 0
3. Laguerre’s differential equation is
xy ′′ + (1 − x)y ′ + λy = 0
Show that when λ = n, a nonnegative integer, there is a polynomial solution Ln (x) (called a Laguerre
polynomial) of degree n with coefficient of xn equal to 1. Determine L0 through L4 .
4. Consider the function y(x) = x2 − 1 defined for x ∈ [0, 4]. Find eight term expansions in terms of a)
Fourier-Sine, b) Fourier-Legendre, c) Fourier-Hermite, d) Fourier-Bessel series and plot your results
on a single graph.
5. Consider the function y(x) = 0, x ∈ [0, 1), y(x) = 2x − 1, x ∈ [1, 2]. Find an eight term Fourier-
Legendre expansion of this function. Plot the function and the eight term expansion for x ∈ [0, 2].
6. Consider the function y(x) = 4x, x ∈ [0, 3]. Find an eight term a) Fourier-Chebyshev and b) Fourier-
sine expansion of this function. Plot the function and the eight term expansions for x ∈ [0, 3]. Which
expansion minimizes the error in representation of the function?
7. Consider the function y(x) = cos2 (x/5). Find an eight term a) Fourier-Laguerre, (x ∈ [0, ∞)), and b)
Fourier-sine (x ∈ [0, 10]) expansion of this function. Plot the function and the eight term expansions
for x ∈ [0, 10]. Which expansion minimizes the error in representation of the function?
Here we will consider what is known as Cartesian index notation as a way to represent vectors
and tensors. In contrast to Chapter 1, which considered general coordinate transformations,
when we restrict our transformations to rotations about the origin, many simplifications
result. For such transformations, the distinction between contravariance and covariance
disappears, as does the necessity for Christoffel symbols, and also the need for an “upstairs-
downstairs” index notation.
Many vector relations can be written in a compact form by using Cartesian index nota-
tion. Let x1 , x2 , x3 represent the three coordinate directions and e1 , e2 , e3 the unit vectors
in those directions. Then a vector u may be written as
u1 X 3
u = u2 = u1 e1 + u2 e2 + u3 e3 =
ui ei = ui ei = ui (6.1)
u3 i=1
where u1 , u2 , and u3 are the three Cartesian components of u. Note that we do not need to
use the summation sign every time if we use the Einstein1 convention to sum from 1 to 3 if
an index is repeated. The single free index on the right side indicating that an ei is assumed.
1
Albert Einstein, 1879-1955, German/American physicist and mathematician.
161
162 CHAPTER 6. VECTORS AND TENSORS
Two additional symbols are needed for later use. They are the Kronecker delta
0 if i 6= j
δij ≡ (6.2)
1 if i = j
The identity
ǫijk ǫlmn = δil δjm δkn + δim δjn δkl + δin δjl δkm − δil δjn δkm − δim δjl δkn − δin δjm δkl (6.4)
relates the two. The following identities are also easily shown:
δii = 3 (6.5)
δij = δji (6.6)
δij δjk = δik (6.7)
ǫijk ǫilm = δjl δkm − δjm δkl (6.8)
ǫijk ǫljk = 2δil (6.9)
ǫijk ǫijk = 6 (6.10)
ǫijk = −ǫikj (6.11)
ǫijk = −ǫjik (6.12)
ǫijk = −ǫkji (6.13)
ǫijk = ǫkij = ǫjki (6.14)
Example 6.1
Let us consider, using generalized coordinates described in an earlier chapter, a trivial identity
transformation from the Cartesian ξ i coordinates to the transformed coordinates xi :
x1 = ξ 1 , x2 = ξ 2 , x3 = ξ 3 .
Here, we are returning to the more general “upstairs-downstairs” index notation of an earlier chapter.
The Jacobian of the transformation is
1 0 0
∂ξ i
J= = 0 1 0 = δji = I.
∂xj
0 0 1
Note that δji has precisely the same properties as δij ; it is zero when i 6= j and unity when i = j. The
metric tensor then is
gij = G = JT · J = I · I = I = δij .
Then we find by the transformation rules that for this transformation, the covariant and contravariant
representations of a general vector u are one and the same:
Consequently, for Cartesian vectors, there is no need to use a notation which distinguishes covariant
and contravariant representations. We will hereafter write all Cartesian vectors with only a subscript
notation.
α ≡ [x1 , x′1 ]
x2
x’2
P
x*2
α α
x’1
x*1’
β β β
α
x*1 x1
π
With β = 2
− α, the angle between the x′1 and x2 axes is
β ≡ [x2 , x′1 ]
The point P is can be represented in both coordinate systems. In the unrotated system, P
is represented by the coordinates:
P : (x∗1 , x∗2 )
In the rotated coordinate system P is represented by
P : (x∗1 ′ , x∗2 ′ )
Trigonometry shows us that
x∗1 ′ = x∗1 cos α + x∗2 cos β (6.15)
x∗1 ′ = x∗1 cos[x1 , x′1 ] + x∗2 cos[x2 , x′1 ] (6.16)
Dropping the stars, and extending to three dimensions, we find that
x′1 = x1 cos[x1 , x′1 ] + x2 cos[x2 , x′1 ] + x3 cos[x3 , x′1 ] (6.17)
Extending to expressions for x′2 and x′3 and writing in matrix form, we get
cos[x1 , x′1 ] cos[x1 , x′2 ] cos[x1 , x′3 ]
( x′1 x′2 x′3 ) = ( x1 x2 x3 ) cos[x2 , x′1 ] cos[x2 , x′2 ] cos[x2 , x′3 ] (6.18)
cos[x3 , x′1 ] cos[x3 , x′2 ] cos[x3 , x′3 ]
Example 6.2
Show for the two-dimensional system described in Figure 6.1 that ℓij ℓkj = δik holds.
Expanding for the two-dimensional system, we get
ℓi1 ℓk1 + ℓi2 ℓk2 = δik .
First, take i = 1, k = 1. We get then
ℓ11 ℓ11 + ℓ12 ℓ12 = δ11 = 1,
cos α cos α + cos(α + π/2) cos(α + π/2) = 1,
cos α cos α + (− sin(α))(− sin(α)) = 1,
cos2 α + sin2 α = 1.
This is obviously true. Next, take i = 1, k = 2. We get then
ℓ11 ℓ21 + ℓ12 ℓ22 = δ12 = 0,
cos α cos(π/2 − α) + cos(α + π/2) cos(α) = 0,
cos α sin α − sin α cos α = 0.
This is obviously true. Next, take i = 2, k = 1. We get then
ℓ21 ℓ11 + ℓ22 ℓ12 = δ21 = 0,
cos(π/2 − α) cos α + cos α cos(π/2 + α) = 0,
sin α cos α + cos α(− sin α) = 0.
Using this, we can easily find the inverse transformation back to the unprimed coordinates
via the following operations:
ℓkj x′j = ℓkj xi ℓij , (6.23)
= ℓij ℓkj xi , (6.24)
= δik xi , (6.25)
= xk , (6.26)
ℓij x′j = xi , (6.27)
xi = ℓij x′j . (6.28)
∂xi
Note that the Jacobian matrix of the transformation is J = ∂x′j
= ℓij . Note it can be shown
that the metric tensor is G = JT · J = ℓjiℓki = δjk = I, so g = 1, and the transformation is
volume preserving. Moreover, since JT · J = I, we see that JT = J−1 .
6.2.1.1 Scalars
A term φ is a scalar if it is invariant under a rotation of coordinate axes.
6.2.1.2 Vectors
A set of three scalars (v1 , v2 , v3 )T is defined as a vector if under a rotation of coordinate axes,
the triple also transforms according to
vj′ = vi ℓij .
A vector associates a scalar with a chosen direction in space by an expression which is
linear in the direction cosines of the chosen direction.
6.2.1.3 Tensors
A set of nine scalars is defined as a second order tensor if under a rotation of coordinate
axes, they transform as
x3
q (3)
Τ 33
Τ32
Τ31
Τ23
(2)
q
Τ13
Τ22
Τ12 Τ21
x2
Τ11
(1)
q
x1
Example 6.3
Returning to generalized coordinate notation, show the equivalence between covariant and con-
travariant representations for pure rotations of a vector v.
Consider then a transformation from a Cartesian space ξ j to a transformed space xi via a pure
rotation:
ξ i = ℓij xj .
Here ℓij is simply a matrix of direction cosines as we have previously defined; we employ the upstairs-
downstairs index notation for consistency. The Jacobian is
∂ξ i
= ℓij .
∂xj
The metric tensor is
∂ξ i ∂ξ i
gkl = = ℓik ℓil = δkl .
∂xk ∂xl
Here we have employed the law of cosines, which is easily extensible to the “upstairs-downstairs”
notation.
Note the vector itself has components that do transform under rotation:
v i = ℓij V j .
Here V j is the contravariant representation of the vector v in the unrotated coordinate system. One
could also show that Vj = V j , as always for a Cartesian system.
qj = ni Tij qT = nT · T (6.29)
Here ni has components which are the direction cosines of the chosen direction. For example
to determine the vector associated with face 2, we choose
0
ni = 1
0
Thus
T11 T12 T13
T
n · T = (0, 1, 0) T21
T22 T23 = (T21 , T22 , T23 ) (6.30)
T31 T32 T33
ni Tij = n1 T1j + n2 T2j + n3 T3j (6.31)
= (0)T1j + (1)T2j + (0)T3j (6.32)
= (T21 , T22 , T23 ) (6.33)
TijT ≡ Tji
so
T11 T21 T31
TijT = T12 T22 T32
T13 T23 T33
A tensor is symmetric if it is equal to its transpose, i.e.
The tensor inner product of a symmetric tensor Sij and anti-symmetric tensor Aij can be
shown to be 0:
Sij Aij = 0
Example 6.4
Show this for a two-dimensional space. Take a general symmetric tensor to be
a b
Sij =
b c
So
Sij Aij = S11 A11 + S12 A12 + S21 A21 + S22 A22 (6.34)
= a(0) + bd − bd + c(0) (6.35)
= 0 (6.36)
1 1 1 1
Tij = Tij + Tij + Tji − Tji (6.37)
|2 {z 2 } |2 {z 2 }
=Tij =0
1 1
= (Tij + Tji ) + (Tij − Tji) (6.38)
|2 {z } |2 {z }
≡T(ij) ≡T[ij]
so with (6.39)
1
T(ij) ≡ (Tij + Tji ) (6.40)
2
1
T[ij] ≡ (Tij − Tji ) (6.41)
2
Tij = T(ij) + T[ij] (6.42)
The first term, T(ij) , is called the symmetric part of Tij ; the second term, T[ij] , is called the
anti-symmetric part of Tij .
1 1 1
di ≡ ǫijk Tjk = ǫijk T(jk) + ǫijk T[jk] (6.43)
2 2 | {z } 2
=0
1
di = ǫijk T[jk] (6.44)
2
Let us find the inverse.
1
ǫilm di = ǫilm ǫijk Tjk (6.45)
2
1
= (δlj δmk − δlk δmj )Tjk (6.46)
2
1
= (Tlm − Tml ) (6.47)
2
= T[lm] (6.48)
T[lm] = ǫilm di (6.49)
T[ij] = ǫkij dk (6.50)
ni Tij = λnj
This defines an eigenvalue problem. Linear algebra gives us the eigenvalues and associated
eigenvectors.
ni Tij = λni δij
ni (Tij − λδij ) = 0
T11 − λ T12 T13
(n1 , n2 , n3 ) T21 T22 − λ T23 = (0, 0, 0)
T31 T32 T33 − λ
We get non-trivial solutions if
T11 − λ T 12 T 13
T21
T22 − λ T23 = 0
T31 T32 T33 − λ
We shall see in later chapters that we are actually finding the so-called left eigenvectors.
These arise with less frequency than the right eigenvectors, which are defined by Tij uj =
λδij uj .
We know from linear algebra that such an equation for a third order matrix gives rise to
a characteristic polynomial for λ of the form
(1) (2) (3)
λ3 − IT λ2 + IT λ − IT = 0, (6.53)
rotated; in constrast, the scalar components Tij will change under rotation. The invariants
can be shown to be given by
(1)
IT = Tii = tr(T), (6.54)
(2) 1 1
IT = (Tii Tjj − Tij Tji ) = (tr(T))2 − tr(T · T) = det(T)tr T−1 , (6.55)
2 2
1
= T(ii) T(jj) + T[ij] T[ij] − T(ij) T(ij) , (6.56)
2
(3)
IT = ǫijk T1i T2j T3k = det (T) . (6.57)
Here “det” denotes the determinant, and “tr” the trace. It can also be shown that if
λ(1) , λ(2) , λ(3) are the three eigenvalues, then the invariants can also be expressed as
(1)
IT = λ(1) + λ(2) + λ(3) , (6.58)
(2) (1) (2) (2) (3) (3) (1)
IT = λ λ +λ λ +λ λ , (6.59)
(3)
IT = λ(1) λ(2) λ(3) . (6.60)
A sketch of a volume element rotated to be aligned with a set of orthogonal principal axes
is shown in Figure 6.3.
If the matrix is not symmetric, the eigenvalues and eigenvectors could be complex. It
is often most physically relevant to decompose a tensor into symmetric and anti-symmetric
parts and find the orthogonal basis vectors and real eigenvalues associated with the sym-
metric part and the dual vector associated with the anti-symmetric part.
In continuum mechanics,
• the symmetric part of a tensor can be associated with deformation along principal axes
Example 6.5
Decompose the tensor given below into a combination of orthogonal basis vectors and a dual vector.
1 1 −2
Tij = 3 2 −3
−4 1 1
x3 x’
3
q(3)
q ( 3’ )
Τ 33
rotate
Τ32
Τ31
Τ23 (2)
q
Τ13
Τ22
Τ12 Τ21
q ( 1’
)
x2
Τ11 x’
1 q (2’ )
(1)
q
x’
2
x1
Figure 6.3: Sketch depicting rotation of volume element to be aligned with principal axes.
Tensor Tij must be symmetric to guarantee existence of orthogonal principal directions.
First
1 2 −3
1
T(ij) = (Tij + Tji ) = 2 2 −1
2
−3 −1 1
0 −1 1
1
T[ij] = (Tij − Tji ) = 1 0 −2
2
−1 2 0
First, get the dual vector di :
1
di = ǫijk T[jk] (6.61)
2
1 1 1
d1 = ǫ1jk T[jk] = (ǫ123 T[23] + ǫ132 T[32] ) = ((1)(−2) + (−1)(2)) = −2 (6.62)
2 2 2
1 1 1
d2 = ǫ2jk T[jk] = (ǫ213 T[13] + ǫ231 T[31] ) = ((−1)(1) + (1)(−1)) = −1 (6.63)
2 2 2
1 1 1
d3 = ǫ3jk T[jk] = (ǫ312 T[12] + ǫ321 T[21] ) = ((1)(−1) + (−1)(1)) = −1 (6.64)
2 2 2
di = (−2, −1, −1)T (6.65)
λ3 − 4λ2 − 9λ + 9 = 0
Example 6.6
For a given tensor, which we will take to be symmetric though the theory applies to non-symmetric
tensors as well,
1 2 4
Tij = T = 2 3 −1 , (6.69)
4 −1 1
(1) (2) (3)
find the three basic tensor invariants, IT , IT , and IT , and show they are truly invariant when the
tensor is subjected to a rotation with direction cosine matrix of
1 q
2√ 1 √
3
6 6
ℓij = L = √1 − √13 √1 (6.70)
3 3
√1 0 − √12
2
We then seek the tensor invariants of T′ . Leaving out some of the details, which are the same as those
for calculating the invariants of the T, we find the invariants indeed are invariant:
(1)
IT = 4.10238 − 0.218951 + 1.11657 = 5, (6.76)
(2) 1 2
IT = (5 − 53) = −14, (6.77)
2
(3)
IT = −66. (6.78)
Finally, we verify that the tensor invariants are indeed related to the principal values (the eigenvalues
of the tensor) as follows
(1)
IT = λ(1) + λ(2) + λ(3) = 5.28675 − 3.67956 + 3.39281 = 5, (6.79)
(2) (1) (2) (2) (3) (3) (1)
IT = λ λ +λ λ +λ λ
= (5.28675)(−3.67956) + (−3.67956)(3.39281) + (3.39281)(5.28675) = −14, (6.80)
(3) (1) (2) (3)
IT = λ λ λ = (5.28675)(−3.67956)(3.39281) = −66. (6.81)
Here the subscript 2 in || · ||2 indicates we are considering a Euclidean norm. In many
sources in the literature this subscript is omitted, and the norm is understood to be the
Euclidean norm. In a more general sense, we can still retain the property of a norm for a
more general p-norm for a three-dimensional vector:
||u||p = (|u1|p + |u2|p + |u3 |p )1/p , 1 ≤ p < ∞.
For example the 1-norm of a vector is the sum of the absolute values of its components:
||u||1 = (|u1 | + |u2| + |u3 |) .
The ∞-norm selects the largest component:
||u||∞ = lim (|u1 |p + |u2 |p + |u3|p )1/p = maxi=1,2,3 |ui|.
p→∞
[u, v, w] = uT · (v × w) (6.84)
= ǫijk ui vj wk (6.85)
Physically it represents the volume of the parallelepiped with edges parallel to the three
vectors.
6.3.5 Identities
Example 6.7
Prove the second identity using Cartesian index notation.
u × (v × w) = ǫijk uj (ǫklm vl wm )
= ǫijk ǫklm uj vl wm
= ǫkij ǫklm uj vl wm
= (δil δjm − δim δjl ) uj vl wm
= uj vi wj − uj vj wi
= uj wj vi − uj vj wi
= (uT · w)v − (uT · v)w
then r(τ ) describes a curve in three-dimensional space. Here τ is a general scalar parameter,
which may or may not have a simple physical interpretation. If we require that the basis
vectors be constants (this is not always the case!), the derivative is
dr(τ )
= r′ (τ ) = x′i (τ )ei = x′i (τ ) (6.91)
dτ
Now r′ (τ ) is a vector that is tangent to the curve. A unit vector in this direction is
r′ (τ )
t= (6.92)
||r′ (τ )||2
where p
||r′(τ )||2 = x′i x′i (6.93)
In the special case in which τ is time t, we denote the derivative by a dot ( ˙ ) notation
rather than a prime (′ ) notation; ṙ is the velocity vector, ẋi its components, and ||ṙ||2 the
magnitude. Note that the unit tangent vector t is not the scalar parameter for time, t. Also
we will occasionally use the scalar components of t: ti , which again are not related to time
t.
If s(t) is the distance along the curve, then
ds2 = dx21 + dx22 + dx23 , (6.94)
q
ds = dx21 + dx22 + dx23 , (6.95)
ds = ||dxi||2 , (6.96)
ds dxi
= (6.97)
dt dt 2
= ||ṙ(t)||2 (6.98)
so that
dr
dt dr dri
t= ds
= , ti = (6.99)
dt
ds ds
Also
Z b Z br Z br
dxi dxi dx1 dx1 dx2 dx2 dx3 dx3
s= ||ṙ(t)||2 dt = dt = + + dt (6.100)
a a dt dt a dt dt dt dt dt dt
is the distance along the curve between t = a and t = b.
Identities:
d du dφ d dui dφ
(φu) = φ + u (φui) = φ + ui
dt dt dt dt dt dt
d T dv duT d dvi dui
(u · v) = uT · + ·v (ui vi ) = ui + vi
dt dt dt dt dt dt
d dv du d dvk duj
(u × v) = u × + ×v (ǫijk uj vk ) = ǫijk uj + ǫijk vk
dt dt dt dt dt dt
CC BY-NC-ND. 28 March 2011, M. Sen, J. M. Powers.
6.4. CALCULUS OF VECTORS 179
Example 6.8
If
r(t) = 2t2 i + t3 j
find the unit tangent at t = 1, and the length of the curve from t = 0 to t = 1.
The derivative is
ṙ(t) = 4ti + 3t2 j
At t = 1,
ṙ(t = 1) = 4i + 3j
so that the unit vector in this direction is
4 3
t= i+ j
5 5
The length of the curve from t = 0 to t = 1 is
Z 1p
s = 16t2 + 9t4 dt
0
1
= (16 + 9t2 )3/2 |10
27
61
=
27
so as ∆θ → 0
t′ − t = −∆θ cos θ i − ∆θ sin θ j
∆t = ∆θ (− cos θ i − sin θ j)
| {z }
unit vector
||∆t||2 = ∆θ
t ’(t + ∆t)
t (t)
∆s
ρ
∆θ
Now for ∆θ → 0,
∆s = ρ∆θ
where ρ is the radius of curvature. So
∆s
||∆t||2 =
ρ
Taking all limits to zero, we get
dt
= 1 (6.101)
ds ρ
2
The term on the right side is often defined as the curvature, κ:
1
κ= .
ρ
ẋi + ẏj
t = (6.104)
(ẋ2
+ ẏ 2 )1/2
i + y ′j
= (6.105)
(1 + (y ′ )2 )1/2
As the second factor of this expression is a unit vector, the preceding scalar is a magnitude.
We define this unit vector to be n, and note that it is orthogonal to the unit tangent vector
t:
−y ′ i + j i + y ′j
nT · t = · , (6.110)
(1 + (y ′)2 )1/2 (1 + (y ′ )2 )1/2
−y ′ + y ′
= , (6.111)
1 + (y ′)2
= 0. (6.112)
dt
Expanding our notion of curvature and radius of curvature, we define ds
such that
dt
= κn, (6.113)
ds
dt
= κ = 1 . (6.114)
ds ρ
2
Thus,
y ′′
κ = , (6.115)
(1 + (y ′)2 )3/2
(1 + (y ′)2 )3/2
ρ = , (6.116)
y ′′
for curves on a plane.
We will first show that t, n, and b form an orthogonal system of unit vectors. We have
already seen that t is a unit vector tangent to the curve. By the product rule for vector
differentiation, we have the identity
dt 1 d T
tT · = (t · t)
ds 2 ds | {z }
=1
1 dt dt dn
= × +t × (6.124)
κ |ds {z ds} ds
=0
dn
= t× (6.125)
ds
db
so we see that ds
is orthogonal to t. In addition, since ||b||2 = 1
db 1 d T
bT · = (b · b)
ds 2 ds
1 d
= (||b||22)
2 ds
= 0
db db
So ds
is orthogonal to b also. Thus ds
must be in the n direction, so that we can write
db
= τn (6.126)
ds
where τ is called the torsion of the curve.
From the relation n = b × t, we get
dn db dt
= ×t+b×
ds ds ds
= τ n × t + b × κn
= −τ b − κt
Summarizing
dt
= κn (6.127)
ds
dn
= −κt − τ b (6.128)
ds
db
= τn (6.129)
ds
These are the Frenet-Serret3 relations. In matrix form, we can say that
t 0 κ 0 t
d
n = −κ 0 −τ n (6.130)
ds
b 0 τ 0 b
Example 6.9
Find the local coordinates, the curvature and the torsion for the helix
dr(t)
= −a sin t i + a cos t j + b k
dt
dr(t) p p
= a2 sin2 t + a2 cos2 t + b2 = a2 + b2
dt
2
This gives us the unit tangent vector t:
dr
dt −a sin t i + a cos t j + b k
t = dr = √
dt 2
a2 + b 2
We also have
s
2 2 2
ds dx dy dz
= + + (6.131)
dt dt dt dt
p
= a2 sin2 t + a2 cos2 t + b2 (6.132)
p
= a2 + b 2 (6.133)
Continuing, we have
dt
dt dt
= ds
ds dt
cos t i + sin t j 1
= −a √ √
2
a +b 2 a + b2
2
= κn
n = −(cos t i + sin t j)
The curvature is
a
κ=
a2 + b2
The radius of curvature is
a2 + b 2
ρ=
a
We also find the unit binormal
b = t×n
i j k
1
= √ −a sin t a cos t b
a2 + b2 − cos t − sin t 0
b sin t i − b cos t j + a k
= √
a2 + b 2
Further identities:
dr d2 r
× 2 = κv 3 b (6.134)
dt dt
dr d r d3 r2 T
× 2
· 3 = −κ2 v 6 τ (6.135)
p dt dt dt
2 2
||r̈||2 ||ṙ||2 − (ṙ · r̈)2
T
= κ (6.136)
||ṙ||32
ds
where v = dt
.
Example 6.10
Find Z
I= uT · dr
C
x
-5
5 -2.5
0
2.5 2.5
y 5
0
-2.5
-5
20
10
Figure 6.5: Three-dimensional curve parameterized by x(t) = a cos t, y(t) = a sin t, z(t) = bt,
with a = 5, b = 1, for t ∈ [0, 25].
if
u = yz i + xy j + xz k
and C goes from (0, 0, 0) to (1, 1, 1) along
(a) the curve x = y 2 = z,
(b) the straight line x = y = z.
We have Z Z
uT · dr = (yz dx + xy dy + xz dz)
C C
(a) Substituting x = y 2 = z, and thus dx = 2ydy, dx = dz we get
Z 1
I = y 3 (2y dy) + y 3 dy + y 4 (2y dy)
0
Z 1
= (2y 4 + y 3 + 2y 5 )dy
0
1
2y 5 y4 y 6
= + +
5 4 3 0
59
=
60
We can achieve the same result in an alternative way that is often more useful for more curves
whose representation is more complicated. Let us paramterize C by taking x = t, y = t2 , z = t. Thus
dx = dt, dy = 2tdt, dz = dt. The end points of C are at t = 0 and t = 1. So the integral is
Z 1
I = (t2 t dt + tt2 (2t) dt + t(t) dt,
0
Z 1
= (t3 + 2t4 + t2 ) dt,
0
1
t4 2t5 t3
= + + ,
4 5 3 0
59
= .
60
(b) Substituting x = y = z =,and thus dx = dy = dz, we get
Z 1 Z 1
I = (x2 dx + x2 dx + x2 dx) = 3x2 dx = x3 |10 = 1.
0 0
Figure 6.6: The vector field u = yzi + xyj + xzk and the curves a) x = y 2 = z; b) x = y = z.
Now we invoke the mean value theorem, which asserts that somewhere within the limits of
integration, the integrand
R takes on its mean value, which we denote with an overline, so
that, for example, V α dV = αV . Thus we get
Z
(∇φ) V = nφ dS, (6.144)
ZS
(∇T · u) V = nT · u dS, (6.145)
S
Z
(∇ × u) V = n × u dS. (6.146)
S
x3
dx 1
dx 3 O x2
x1 dx 2
where φ(r) is a scalar field, and u(r) is a vector field. V is the region enclosed within a
closed surface S, and n is the unit normal to an element of the surface dS. Here “grad” is
the gradient operator, “div” is the divergence operator, and “curl” is the curl operator.
Consider the element of volume in Cartesian coordinates shown in Figure 6.7. The
differential operations in this coordinate system can be deduced from the definitions and
written in terms of the vector operator ∇:
∂
∂x1
∂ ∂ ∂ ∂
∇ = e1 + e2 + e3 = ∂x∂ 2 = (6.150)
∂x1 ∂x2 ∂x3 ∂ ∂xi
∂x3
Note, for full clarity, we adopt the unconventional, but correct row vector operator
∇T = ( ∂x∂ 1 ∂
∂x2
∂
∂x3
) (6.151)
grad (φ(xi ))
We take the reference value of φ to be at the origin O. Consider first the x1 variation. At
O, x1 = 0, and our function takes the value of φ. At the faces a distance x1 = ± dx21 away
from O in the x1 -direction, our function takes a value of
∂φ dx1
φ±
∂x1 2
∂φ ∂φ ∂φ
= e1 + e2 + e3
∂x1 ∂x2 ∂x3
∂φ ∂φ
= ei =
∂xi ∂xi
= ∇φ (6.152)
The derivative of φ on a particular path is called the directional derivative. If the path
has a unit tangent t , the derivative in this direction is
∂φ
(∇φ)T · t = ti (6.153)
∂xi
∂φ
dφ = dxi
∂xi
= (∇φ)T · dr
Since dr is tangent to the surface, ∇φ must be normal to it. The tangent plane at r = r0 is
defined by the position vector r such that
(∇φ)T · (r − r0 ) = 0 (6.154)
Example 6.11
At the point (1,1,1), find the unit normal to the surface
z 3 + xz = x2 + y 2
Define
φ(x, y, z) = z 3 + xz − x2 − y 2 = 0
-2 -1 0 1 2
4
z
2
-2
-1
0
1
x
2
A normal at (1,1,1) is
6.6.2 Divergence
6.6.2.1 Vectors
Equation (6.148) becomes
1 ∂u1 dx1 ∂u1 dx1
div u = lim u1 + dx2 dx3 − u1 − dx2 dx3
V →0 V ∂x1 2 ∂x1 2
+ similar terms from the x2 and x3 faces
6.6.2.2 Tensors
The extension to tensors is straightforward
divT = ∇T · T (6.156)
∂Tij
= (6.157)
∂xi
Two similar terms appear on the opposite face, whose unit vector points in the −e1 direction.
Carrying out the integration then for equation (6.149), one gets
1 ∂u2 dx1 ∂u3 dx1
curl u = lim u2 + e3 dx2 dx3 − u3 + e2 dx2 dx3
V →0 V ∂x1 2 ∂x1 2
∂u2 dx1 ∂u3 dx1
− u2 − e3 dx2 dx3 + u3 − e2 dx2 dx3
∂x1 2 ∂x1 2
+ similar terms from the x2 and x3 faces
e1 e2 e3
= ∂x∂ 1 ∂x∂ 2 ∂
∂x3
u1 u2 u3
∂uk
= ǫijk
∂xj
= ∇×u (6.162)
6.6.4 Laplacian
6.6.4.1 Scalar
The Laplacian4 is simply div grad φ, and can be written as
∂2φ
∇T · (∇φ) = ∇2 φ = (6.163)
∂xi ∂xi
6.6.4.2 Vector
One of the identities below
∇2 u = ∇T · ∇u = ∇(∇T · u) − ∇ × (∇ × u) (6.164)
6.6.5 Identities
∇ × (∇φ) = 0 (6.165)
∇T · (∇ × u) = 0 (6.166)
∇T · (φu) = φ∇T · u + (∇φ)T · u (6.167)
∇ × (φu) = φ∇ × u + ∇φ × u (6.168)
4
Pierre-Simon Laplace, 1749-1827, Norman-born French mathematician.
∇T · (u × v) = vT · (∇ × u) − uT · (∇ × v) (6.169)
∇ × (u × v) = (vT · ∇)u − (uT · ∇)v + u(∇T · v) − v(∇T · u) (6.170)
∇(uT · v) = (uT · ∇)v + (vT · ∇)u + u × (∇ × v) + v × (∇ × u) (6.171)
∇ × (∇ × u) = ∇(∇T · u) − ∇T · ∇u (6.172)
Example 6.12
Show that
∇ · ∇T u = ∇(∇T · u) − ∇ × (∇ × u)
Going from right to left
∂ ∂uj ∂ ∂um
∇(∇T · u) − ∇ × (∇ × u) = − ǫijk ǫklm
∂xi ∂xj ∂xj ∂xl
∂ ∂uj ∂ ∂um
= − ǫkij ǫklm
∂xi ∂xj ∂xj ∂xl
∂ 2 uj ∂ 2 um
= − (δil δjm − δim δjl )
∂xi ∂xj ∂xj ∂xl
2 2 2
∂ uj ∂ uj ∂ ui
= − +
∂xi ∂xj ∂xj ∂xi ∂xj ∂xj
∂ ∂ui
=
∂xj ∂xj
= ∇T · ∇u
C2
C3
C1 x
0 1
Figure 6.9: Sketch of vector field u = yi + 2xyj and closed contour integral C.
Example 6.13
Show that Green’s theorem is valid if u = y i + 2xy j, and C consists of the straight lines (0,0) to
(1,0) to (1,1) to (0,0).
I Z Z Z
uT · dr = uT · dr + uT · dr + uT · dr
C C1 C2 C3
where C1 , C2 , and C3 are the straight lines (0,0) to (1,0), (1,0) to (1,1), and (1,1) to (0,0), respectively.
This is sketched in Figure 6.9. For this problem we have
C1 : y = 0, dy = 0, x ∈ [0, 1], u = 0
C2 : x = 1, dx = 0, y ∈ [0, 1], u = y i + 2y j
C3 : x = y, dx = dy, x ∈ [1, 0], y ∈ [1, 0] u = x i + 2x2 j
Thus
I Z 1 Z 1 Z 0
u · dr = (0 i + 0 j) · (dx i) + (y i + 2y j) · (dy j) + (x i + 2x2 j) · (dx i + dx j)
C
|0 {z } |0 {z } |1 {z }
C1 C2 C3
Z 1 Z 0
= 2y dy + (x + 2x2 ) dx
0 1
0
1 1 2 2 3 1 2
= y 2 0 + x + x =1− −
2 3 1 2 3
1
= −
6
On the other hand
ZZ Z 1 Z x
∂uy ∂ux
− dx dy = (2y − 1) dy dx
D ∂x ∂y 0 0
Z 1 x
= y 2 − y 0 dx
0
Z 1
= (x2 − x) dx
0
1
x3 x2
= −
3 2 0
1 1 1
= − =−
3 2 6
where dV an element of volume, dS is an element of the surface, and n (or ni ) is the outward
unit normal to it. The divergence theorem is also known as Gauss’s theorem. It extends to
tensors of arbitrary order: Z Z
∂Tijk...
Tijk...ni dS = dV
S V ∂xi
Note if Tijk... = C then we get Z
ni dS = 0
S
-1 0
1
1.0
0.5
z 0.0
-0.5
-1.0
-1 0 1
Z b
dφ
φ(b) − φ(a) = dx (6.176)
a dx
Here the end points play the role of the surface integral, and the integral on x plays the
role of the volume integral.
Example 6.14
Show that the divergence theorem is valid if
u = x i + y j + 0k
and S is the closed surface which consists of a circular base and the hemisphere of unit radius with
center at the origin and z ≥ 0, that is,
x2 + y 2 + z 2 = 1
r = 1.
A sketch of the surface of interest along with the vector field is shown in Figure 6.10.
Thus on H
Interchanging φ and ψ in the above and subtracting, we get Green’s second identity
Z Z
T
(φ∇ψ − ψ∇φ) · n dS = (φ∇2 ψ − ψ∇2 φ) dV (6.181)
S V
Z Z
∂ψ ∂φ ∂2ψ ∂2φ
φ −ψ ni dS = φ −ψ dV (6.182)
S ∂xi ∂xi V ∂xi ∂xi ∂xi ∂xi
where n is the unit vector normal to the element dS, and dr an element of curve C.
Example 6.15
Evaluate Z
I= (∇ × u)T · n dS
S
5
George Gabriel Stokes, 1819-1903, Irish-born English mathematician.
y 0
-2
2 z
1
-1
0 0
1
x
C can be represented by the parametric equations x = cos t, y = 2 sin t. Thus dy = 2 cos t dt, so that
Z 2π
I = 2 cos4 t dt
0
Z 2π
1 1 3
= 2 cos 4t + cos 2t + dt
0 8 2 8
2π
1 1 3
= 2 sin 4t + sin 2t + t
32 4 8 0
3
= π
2
A sketch of the surface of interest along with the vector field is shown in Figure 6.11. The curve C is
on the boundary z = 0.
Z Z Z
d ∂Tjk... (xi , t)
Tjk...(xi , t) dV = dV + nm wm Tjk....(xi , t) dS (6.185)
dt V (t) V (t) ∂t S(t)
Here the volume changes due to the net surface motion. In one dimension Tjk...(xi , t) = f (x, t)
we get
Z x=b(t) Z x=b(t)
d ∂f db da
f (x, t) dx = dx + f (b(t), t) − f (a(t), t) (6.188)
dt x=a(t) x=a(t) ∂t dt dt
where s 2 2 2
∂x1 ∂x2 ∂x3
hi = + + (6.190)
∂qi ∂qi ∂qi
We can show that
1 ∂φ 1 ∂φ 1 ∂φ
grad φ = e1 + e2 + e3
h1 ∂q1 h2 ∂q2 h3 ∂q3
1 ∂ ∂ ∂
div u = (u1 h2 h3 ) + (u2 h3 h1 ) + (u3 h1 h2 )
h1 h2 h3 ∂q1 ∂q2 ∂q3
h1 e1 h2 e2 h3 e3
1 ∂ ∂ ∂
curl u = ∂q ∂q ∂q
h1 h2 h3 1 2 3
u1 h1 u2 h2 u3 h3
1 ∂ h2 h3 ∂φ ∂ h3 h1 ∂φ ∂ h1 h2 ∂φ
div grad φ = + +
h1 h2 h3 ∂q1 h1 ∂q1 ∂q2 h2 ∂q2 ∂q3 h3 ∂q3
Example 6.16
Find expressions for the gradient, divergence, and curl in cylindrical coordinates (r, θ, z) where
x1 = r cos θ
x2 = r sin θ
x3 = z
The 1,2 and 3 directions are associated with r, θ, and z, respectively. From equation (6.190) the scale
factors are
s 2 2 2
∂x1 ∂x2 ∂x3
hr = + +
∂r ∂r ∂r
p
2
= cos2 θ + sin θ
= s 1
2 2 2
∂x1 ∂x2 ∂x3
hθ = + +
∂θ ∂θ ∂θ
p
2
= r2 sin θ + r2 cos2 θ
= s r
2 2 2
∂x1 ∂x2 ∂x3
hz = + +
∂z ∂z ∂z
= 1
so that
∂φ 1 ∂φ ∂φ
grad φ = er + eθ + ez
∂r r ∂θ ∂z
1 ∂ ∂ ∂ ∂ur ur 1 ∂uθ ∂uz
div u = (ur r) + (uθ ) + (uz r) = + + +
r ∂r ∂θ ∂z ∂r r r ∂θ ∂z
e reθ ez
1 ∂r ∂ ∂
curl u =
r ∂r ∂θ ∂z
ur uθ r uz
Problems
1. Find the angle between the planes
3x − y + 2z = 2
x − 2y = 1
2. Find the curve of intersection of the cylinders x2 + y 2 = 1 and y 2 + z 2 = 1. Determine also the radius
of curvature of this curve at the points (0,1,0) and (1,0,1).
3. Show that for a curve r(t)
dt d2 t
tT · × = κ2 τ
ds ds2
drT d2 r d3 r
ds · ds 2 × ds3
d2 rT d2 r
= τ
ds2 · ds2
where t is the unit tangent, s is the length along the curve, κ is the curvature, and τ is the torsion.
4. Find the equation for the tangent to the curve of intersection of x = 2 and y = 1 + xz sin y 2 z at the
point (2, 1, π).
5. Find the curvature and torsion of the curve r(t) = 2ti + t2 j + 2t3 k at the point (2, 1, 2).
6. Apply Stokes’s theorem to the plane vector field u(x, y) = uxHi + uy j and a closed curve enclosing a
plane region. What is the result called? Use this result to find C uT · dr, where u = −yi + xj and the
integration is counterclockwise along the sides C of the trapezoid with corners at (0,0), (2,0), (2,1),
and (1,1).
7. Orthogonal bipolar coordinates (u, v, w) are defined by
α sinh v
x =
cosh v − cos u
α sin u
y =
cosh v − cos u
z = w
For α = 1, plot some of the surfaces of constant x and y in the u − v plane.
8. Using Cartesian index notation, show that
∇ × (u × v) = (vT · ∇)u − (uT · ∇)v + u(∇T · v) − v(∇T · u)
where u and v are vector fields.
9. Consider two Cartesian
√ coordinate systems:
√ S with unit vectors (i, j, k), and S ′ with (i′ , j′ , k′ ), where
′ ′ ′
i = i, j = (j − k)/ 2, k = (j + k)/ 2. The tensor T has the following components in S:
1 0 0
0 −1 0
0 0 2
Find its components in S ′ .
10. Find the matrix A that operates on any vector of unit length in the x-y plane and turns it through
an angle θ around the z-axis without changing its length. Show that A is orthogonal; that is that all
of its columns are mutually orthogonal vectors of unit magnitude.
11. What is the unit vector normal to the plane passing through the points (1,0,0), (0,1,0) and (0,0,2)?
12. Prove the following identities using Cartesian index notation:
(a) (a × b)T · c = aT · (b × c)
(b) a × (b × c) = b(aT · c) − c(aT · b)
T
(c) (a × b)T · (c × d) = ((a × b) × c) · d
13. The position of a point is given by r = ia cos ωt + jb sin ωt. Show that the path of the point is an
ellipse. Find its velocity v and show that r × v = constant. Show also that the acceleration of the
point is directed towards the origin and its magnitude is proportional to the distance from the origin.
14. System S is defined by the unit vectors e1 , e2 , and e3 . Another Cartesian system S ′ is defined by
unit vectors e′1 , e′2 , and e′3 in directions a, b, and c where
a = e1
b = e2 − e3
(a) Find e′1 , e′2 , e′3 , (b) find the transformation array Aij , (c) show that δij = Aki Akj is satisfied, and
(d) find the components of the vector e1 + e2 + e3 in S ′ .
H
15. Use Green’s theorem to calculate C uT · dr, where u = x2 i + 2xyj, and C is the counterclockwise
path around a rectangle with vertices at (0,0), (2,0), (0,4) and (2,4).
16. Derive an expression for the gradient, divergence, curl, and Laplacian operators in orthogonal paraboloidal
coordinates
x = uv cos θ
y = uv sin θ
1 2
z = (u − v 2 )
2
Determine the scale factors. Find ∇φ, ∇T · u, ∇ × u, and ∇2 φ in this coordinate system.
17. Derive an expression for the gradient, divergence, curl and Laplacian operators in orthogonal parabolic
cylindrical coordinates (u, v, w) where
x = uv
1 2
y = (u − v 2 )
2
z = w
x = a cosh u cos v
y = a sinh u sin v
z = z
where u ∈ [0, ∞), v ∈ [0, 2π) and z ∈ (−∞, ∞). Determine ∇f, ∇T · u, ∇ × u and ∇2 f in this system,
where f is a scalar field and u is a vector field.
19. Determine a unit vector in the plane of the vectors i − j and j + k and perpendicular to the vector
i − j + k.
20. Determine a unit vector perpendicular to the plane of the vectors a = i + 2j − k, b = 2i + j + 0k.
21. Find the curvature and the radius of curvature of y = a sin x at the peaks and valleys.
22. Determine the unit vector normal to the surface x3 − 2xyz + z 3 = 0 at the point (1,1,1).
23. Show using indicial notation that
∇ × ∇φ = =0
∇T · ∇ × u = 0
∇(uT · v) = (uT · ∇)v + (vT · ∇)u + u × (∇ × v) + v × (∇ × u)
1
∇(uT · u) = (uT · ∇)u + u × (∇ × u)
2
∇T · (u × v) = vT · ∇ × u − uT · ∇ × v
∇ × (∇ × u) = ∇(∇T · u) − ∇2 u
∇ × (u × v) = (vT · ∇)u − (uT · ∇)v + u(∇T · v) − v(∇T · u)
∂2
24. Show that the Laplacian operator ∂xi ∂xi has the same form in S and S ′ .
25. If
x1 x22 2x3 x1 − x2
Tij = x2 x1 x1 x3 x33 + 1
0 2 2x2 − x3
a) Evaluate Tij at P : (1, 2, 3)
b) find T(ij) and T[ij] at P
c) find the associated dual vector di
d) find the principal values and the orientations of each associated normal vector for the symmetric
part of Tij evaluated at P
e) evaluate the divergence of Tij at P
f) evaluate the curl of the divergence of Tij at P
26. Consider the tensor
2 −1 2
Tij = 3 1 0
0 1 3
defined in a Cartesian coordinate system. Consider the vector associated with the plane whose normal
points in the direction (2, 3, −1). What is the magnitude of the component of the associated vector
that is aligned with the normal to the plane?
Linear analysis
7.1 Sets
Consider two sets A and B. We use the following notation
x∈A x is an element of A
x∈
/A x is not an element of A
A=B A and B have the same elements
A⊂B the elements of A also belong to B
A∪B set of elements that belong to A or B
A∩B set of elements that belong to A and B
A−B set of elements that belong to A but not to B
If A ⊂ B, then B − A is the complement of A in B.
Some sets that are commonly used are:
Z set of all integers
N set of all positive integers
Q set of all rational numbers
R set of all real numbers
R+ set of all non-negative real numbers
207
208 CHAPTER 7. LINEAR ANALYSIS
• An open interval (a, b) does not include the end points, so that if x ∈ (a, b), then
a < x < b. In set notation this is {x ∈ R : a < x < b} if x is real
• A closed interval [a, b] includes the end points. If x ∈ [a, b], then a ≤ x ≤ b. In set
notation this is {x ∈ R : a ≤ x ≤ b} if x is real
• A set A ⊂ R is bounded from above if there exists a real number, called the upper
bound, such that for every x ∈ A is less than or equal to that number.
• The least upper bound or supremum is the minimum of all upper bounds.
• In a similar fashion, a set A ⊂ R can be bounded from below, in which case it will
have a greatest lower bound or infimum.
• A set which has no elements is the empty set {}, also known as the null set ∅. Note
the set with 0 as the only element, 0, is not empty.
• A set that is either finite, or for which each element can be associated with a member
of N is said to be countable. Otherwise the set is uncountable.
f(t)
t0 t1 t2 t tk t t
ξ
1
ξ
2
k-1 ξ n-1 ξ n n t
k
a b
If I exists and is independent of the manner of subdivision, then f (t) is Riemann2 integrable
in [a, b].
The Riemann integration process is sketched in Figure 7.1.
Example 7.1
Determine if the function f (t) Riemann integrable in [0, 1] where
0 if t is rational
f (t) =
1 if t is irrational
2
Georg Friedrich Bernhard Riemann, 1826-1866, Hanover-born German mathematician.
f(t)
y
n
y
n-1
y
k
y
k-1
y
1
y0
e1 ek en t
a b
ek = {t : yk−1 ≤ f (t) ≤ yk }
(ek is the set of all t’s for which f (t) is bounded between two values, yk−1 and yk ). Also let
the sum In be defined as
In = y1 m(e1 ) + y2 m(e2 ) + · · · + yn m(en ) (7.2)
Example 7.2
To integrate the function in the previous example, we observe first that the set of rational and
irrational numbers in [0,1] has measure zero and 1 respectively. Thus from equation (7.2) the Lebesgue
3
Henri Lèon Lebesgue, 1875-1941, French mathematician.
integral exists, and is equal to 1. Loosely speaking, the reason is that the rationals are not dense in
[0, 1] while the irrationals are dense in [0, 1]. That is to say every rational number exists in isolation
from other rational numbers and surrounded by irrationals. Thus the rationals exist as isolated points
on the real line; these points have measure 0; The irrationals have measure 1 over the same interval;
hence the integral is In = y1 m(e1 ) + y2 m(e2 ) = 1(1) + 0(0) = 1.
The Riemann integral is based on the concept of the length of an interval, and the
Lebesgue integral on the measure of a set. When both integrals exist, their values are the
same. If the Riemann integral exists, the Lebesgue integral also exists. The converse is not
necessarily true.
The importance of the distinction is subtle. It can be shown that certain integral oper-
ators which operate on Lebesgue integrable functions are guaranteed to generate a function
which is also Lebesgue integrable. In contrast, certain operators operating on functions which
are at most Riemann integrable can generate functions which are not Riemann integrable.
6. ∀ a, b ∈ F, ∀x ∈ S, (a + b)x = ax + bx. For all a and b which are in F and for all x
which are in S, the addition operator distributes onto multiplication.
7. ∀ a ∈ F, ∀ x, y ∈ S, a(x + y) = ax + ay
8. ∀ a, b ∈ F, ∀ x ∈ S, a(bx) = (ab)x
Such a set is called a linear space or vector space over the field F, and its elements are
called vectors. We will see that our definition is inclusive enough to include elements which
are traditionally thought of as vectors (in the sense of a directed line segment), and some
which are outside of this tradition. Note that typical vector elements x and y are no longer
indicated in bold. However, they are in general not scalars, though in special cases, they can
be.
The element 0 ∈ S is called the null vector. Examples of vector spaces S over the field of
real numbers (i.e. F : R) are:
1. S : R1 . Set of real numbers, x = x1 , with addition and scalar multiplication defined as
usual; also known as S : R.
2. S : R2 . Set of ordered pairs of real numbers, x = (x1 , x2 )T , with addition and scalar
multiplication defined as:
x1 + y1
x+y = = (x1 + y1 , x2 + y2 )T ,
x2 + y2
αx1
αx = = (αx1 , αx2 )T ,
αx2
where
x1 T 2 y1
x= = (x1 , x2 ) ∈ R , y= = (y1 , y2 )T ∈ R2 , α ∈ R1 .
x2 y2
4. S : R∞ . Set of an infinite number of real numbers, x = (x1 , x2 , · · ·)T , with addition and
scalar multiplication defined similar to the above. Functions, e.g. x = 3t2 + t, t ∈ R1
generate vectors x ∈ R∞ .
10. S : C[a, b] Set of real-valued continuous functions, x(t) for t ∈ [a, b] ∈ R1 with addition
and scalar multiplication defined as usual.
11. S : C n [a, b] Set of real-valued functions x(t) for t ∈ [a, b] with continuous nth derivative
with addition and scalar multiplication defined as usual; n ∈ N.
12. S : L2 [a, b] Set of real-valued functions x(t) such that x(t)2 is Lebesgue integrable in
t ∈ [a, b] ∈ R1 , a < b, with addition and multiplication by a scalar defined as usual.
Note that the integral must be finite.
13. S : Lp [a, b] Set of real-valued functions x(t) such that |x(t)|p , p ∈ [1, ∞), is Lebesgue
integrable in t ∈ [a, b] ∈ R1 , a < b, with addition and multiplication by a scalar defined
as usual. Note that the integral must be finite.
14. S : Lp [a, b] Set of complex-valued functions x(t) such that |x(t)|p , p ∈ [1, ∞) ∈ R1 , is
Lebesgue integrable in t ∈ [a, b] ∈ R1 , a < b, with addition and multiplication by a
scalar defined as usual.
1 2
Pn ∂u 2
15. S : W2 (G), Set of real-valued functions u(x) such that u(x) and i=1 ∂xi are
Lebesgue integrable in G, where x ∈ G ∈ Rn , n ∈ N. This is an example of a Sobolov5
space, which is useful in variational calculus and the finite element method. Sobolov
space W21 (G) is to Lebesgue space L2 [a, b] as the real space R1 is to the rational space
Q1 . That is Sobolov space allows a broader class of functions to be solutions to physical
problems. See Zeidler.
m p
Pn ∂u p
16. S : Wp (G), Set of real-valued functions u(x) such that |u(x)| and i=1 ∂xi +
p
m
· · · ∂∂xmu is Lebesgue integrable in G, where x ∈ G ∈ Rn , n, m ∈ N (may be in error
i
due to neglect of mixed partial derivatives!).
17. S : Pn Set of all polynomials of degree ≤ n with addition and multiplication by a scalar
defined as usual; n ∈ N.
Some examples of sets that are not vector spaces are Z and N over the field R for the same
reason that they do not form a field, namely that they are not closed over the multiplication
operation.
5
Sergei Lvovich Sobolev, 1908-1989, St. Petersburg-born Russian physicist and mathematician.
S = S1 ⊕ S2 (7.4)
• The set of all linear combination of k vectors {x1 , x2 , · · · , xk } of a vector space constitute
a subspace of the vector space.
• If the vector space S contains a set of n linearly independent set of vectors, and any
set with (n + 1) elements is linearly dependent, then the space is said to be finite
dimensional, and n is the dimension of the space. If n does not exist, the space is
infinite dimensional.
• A set of vectors in a linear space S is convex iff ∀x, y ∈ S and α ∈ [0, 1] ∈ R1 implies
αx + (1 − α)y ∈ S. For example if we consider S to be a subspace of R2 , that is a
region of the x, y plane, S is convex if for any two points in S, all points on the line
segment between them also lie in S. Spaces with lobes are not convex. Functions f
are convex iff the space on which they operate are convex and if f (αx + (1 − α)y) ≤
αf (x) + (1 − α)f (y) ∀ x, y ∈ S, α ∈ [0, 1] ∈ R1 .
1. ||x|| ≥ 0.
The norm is a natural generalization of the length of a vector. All properties of a norm can
be cast in terms of ordinary finite dimensional Euclidean vectors, and thus have geometrical
interpretations. The first property says length is greater than or equal to zero. The second
says the only vector with zero length is the zero vector. The third says the length of a scalar
multiple of a vector is equal to the magnitude of the scaler times the length of the original
vector. The Minkowski inequality is easily understood in terms of vector addition. If we add
vectorially two vectors x and y, we will get a third vector whose length is less than or equal
to the sum of the lengths of the original two vectors. We will get equality when x and y
point in the same direction. The interesting generalization is that these properties hold for
the norms of functions as well as ordinary geometric vectors.
Examples of norms are:
1. x ∈ R1 , ||x|| = |x|. This space is also written as ℓ1 (R1 ) or in abbreviated form ℓ11 . The
subscript on ℓ in either case denotes the type of norm; the superscript in the second
form denotes the dimension of the space. Another way to denote this norm is ||x||1 .
p √
2. x ∈ R2 , x = (x1 , x2 )T , the Euclidean norm ||x|| = ||x||2 = + x21 + x22 = + xT x. We
can call this normed space E2 , or ℓ2 (R2 ), or ℓ22 .
p √
3. x ∈ Rn , x = (x1 , x2 , · · · , xn )T , ||x|| = ||x||2 = + x21 + x22 + · · · + x2n = + xT x. We
can call this norm the Euclidean norm and the normed space Euclidean En , or ℓ2 (Rn )
or ℓn2 .
qR
b
10. x ∈ L2 [a, b], ||x|| = ||x||2 = + a
x(t)2 dt; t ∈ [a, b] ∈ R1 .
R 1/p
b
11. x ∈ Lp [a, b], ||x|| = ||x||p = + a
|x(t)|p dt ; t ∈ [a, b] ∈ R1 .
qR qR
b b
12. x ∈ L2 [a, b], ||x|| = ||x||2 = + a
|x(t)|2 dt = + a
x(t)x(t) dt; t ∈ [a, b] ∈ R1 .
• The metric or distance between x and y is defined by d(x, y) = ||x − y||. This a natural
metric induced by the norm. Thus ||x|| is the distance between x and the null vector.
• The diameter of a set of vectors is the supremum (i.e. least upper bound) of the
distance between any two vectors of the set.
• The sequence x(1) , x(2) , · · · ∈ S, where S is a normed vector space, converges if there
exists an x ∈ S such that limn→∞ ||x(n) − x|| = 0. Then x is the limit point of the
sequence, and we write limn→∞ x(n) = x or x(n) → x.
• Every convergent sequence is a Cauchy sequence, but the converse is not true.
7
Augustin-Louis Cauchy, 1789-1857, French mathematician and physicist.
• Norms || · ||i and || · ||j in S are equivalent if there exist a, b > 0 such that, for any x ∈ S,
• In a finite-dimensional vector space, any norm is equivalent to any other norm. So,
the convergence of a sequence in such a space does not depend on the choice of norm.
• (z1 + z2 ) = z1 + z2
• (z1 − z2 ) = z1 − z2
• z1 z2 = z1 z2
• zz21 = zz12
We also recall that the modulus of z, |z| has the following properties:
Example 7.3
Consider x ∈ R3 and take
1
x = −4 .
2
Find the norm if x ∈ ℓ31 (absolute value norm), x ∈ ℓ32 (Euclidean norm), if x = ℓ33 (another norm), and
if x ∈ ℓ3∞ (maximum norm).
By the definition of the absolute value norm for x ∈ ℓ31 ,
we get
||x||1 = |1| + | − 4| + |2| = 1 + 4 + 2 = 7.
8
Stefan Banach, 1892-1945, Polish mathematician.
Now consider the Euclidean norm for x ∈ ℓ32 . By the definition of the Euclidean norm,
q
||x|| = ||x||2 + x21 + x22 + x23 ,
we get p √ √
||x||2 = + 12 + (−4)2 + 22 = 1 + 16 + 4 = + 21 ∼ 4.583.
Since the norm is Euclidean, this is the ordinary length of the vector.
For the norm, x ∈ ℓ33 , we have
1/3
||x|| = ||x||3 = + |x1 |3 + |x2 |3 + |x3 |3 ,
so 1/3 1/3
||x||3 = + |1|3 + | − 4|3 + |2|3 = (1 + 64 + 8) ∼ 4.179
so
1/p
||x||∞ = lim + (|1|p + | − 4|p + |2|p ) = 4.
p→∞
This picks out the magnitude of the component of x whose magnitude is maximum.
Note that as p increases the norm of the vector decreases.
Example 7.4
For x ∈ ℓ2 (C2 ), find the norm of
i 0 + 1i
x= = .
1 1 + 0i
so s
0 + 1i
||x||2 = + ( 0 + 1i 1 + 0i ) ,
1 + 0i
q p
||x||2 = + (0 + 1i)(0 + 1i) + (1 + 0i)(1 + 0i) = + (0 − 1i)(0 + 1i) + (1 − 0i)(1 + 0i),
p √
||x||2 = + −i2 + 1 = + 2.
√
Note that if we were negligent in the use of the conjugate and defined the norm as ||x||2 = + xT x,
we would obtain
s
√ i p √
T
||x||2 = + x x = + ( i 1 ) = + i2 + 1 = + −1 + 1 = 0!
1
Example 7.5
Consider x ∈ L2 [0, 1] where x(t) = 2t; t ∈ [0, 1] ∈ R1 . Find ||x||.
By the definition of the norm for this space, we have
s
Z 1
||x|| = ||x||2 = + x2 (t) dt
0
Z 1 Z 1 Z 1 3 1
t
||x||22 = x(t)x(t) dt = (2t)(2t) dt = 4 t2 dt = 4 ,
0 0 0 3 0
3
1 03 4
||x||22 = 4 − = ,
3 3 3
√
2 3
||x||2 = ∼ 1.1547.
3
Example 7.6
Consider x ∈ L3 [−2, 3] where x(t) = 1 + 2it; t ∈ [−2, 3] ∈ R1 . Find ||x||.
By the definition of the norm we have
Z 3 1/3
||x|| = ||x||3 = + |1 + 2it|3 dt
−2
Z 3 3/2 1/3
||x||3 = + (1 + 2it) (1 + 2it) dt
−2
Z 3 3/2
||x||33 = (1 + 2it) (1 + 2it) dt
−2
Z 3
||x||33 = ((1 − 2it) (1 + 2it))3/2 dt
−2
Z 3
3/2
||x||33 = 1 + 4t2 dt
−2
p 3
5t 3
||x||33 = 1 + 4t2 + t3 + sinh−1 (2t)
8 16 −2
√ −1
37 17 3 sinh (4) 3 √
||x||33 = + + 154 17 + sinh−1 (6) ∼ 214.638
4 16 16
||x||3 ∼ 5.98737
Example 7.7
Consider x ∈ Lp [a, b] where x(t) = c; t ∈ [a, b] ∈ R1 , c ∈ C1 . Find ||x||. Let us take the complex
constant c = α + iβ, α ∈ R1 , β ∈ R1 . Then
1/2
|c| = α2 + β 2 .
Now
Z !1/p
b
||x|| = ||x||p = |x(t)|p dt
a
Z !1/p
b
2 2 p/2
||x||p = α +β dt
a
Z !1/p
b
2 2 p/2
||x||p = α +β dt
a
p/2 1/p
||x||p = α2 + β 2 (b − a)
1/2
||x||p = α2 + β 2 (b − a)1/p
||x||p = |c|(b − a)1/p
Note the norm is proportional to the magnitude of the complex constant c. For finite p, it is also
increases with the extent of the domain b − a. For infinite p, it is independent of the length of the
domain, and simply selects the value |c|. This is consistent with the norm in L∞ selecting the maximum
value of the function.
Example 7.8
Consider x ∈ Lp [0, b] where x(t) = 2t2 ; t ∈ [0, b] ∈ R1 . Find ||x||. Now
Z !1/p
b
||x|| = ||x||p = |x(t)|p dt
0
Z !1/p
b
2 p
||x||p = |2t | dt
0
Z !1/p
b
p 2p
||x||p = 2 t dt
0
b !1/p
2p t2p+1
||x||p =
2p + 1 0
p 2p+1 1/p
2 b
||x||p =
2p + 1
2p+1
2b p
||x||p =
(2p + 1)1/p
1/p 2p+1
Note as p → ∞ that (2p + 1) → 1, and p → 2, so
Example 7.9
Consider u ∈ W12 (G) with u(x) = 2x4 ; x ∈ [0, 3] ∈ R1 . Find ||u||.
We note that here that n = 1; the consequent one-dimensional domain G = [0, 3] is a closed interval
on the real number line. For more general problems, it can be areas, volumes, or n-dimensional regions
of space. Also here m = 1 and p = 2, so we require u ∈ L2 [0, 3] and ∂u ∂x ∈ L2 [0, 3], which for our choice
of u, is satisfied. The formula for the norm in W21 [0, 3] is
s
Z 3
du du
||u|| = ||u||1,2 = + u(x)u(x) + dx,
0 dx dx
s
Z 3
||u||1,2 = + ((2x4 )(2x4 ) + (8x3 )(8x3 )) dx,
0
s
Z 3
||u||1,2 = + (4x8 + 64x6 ) dx,
0
s 3 r
4x9 64x7 69
||u||1,2 = + + = 54 7 ∼ 169.539,
9 7 0
Example 7.10
Consider the sequence of vectors {x(1) , x(2) , . . .} ∈ Q3 , where Q3 is the space of rational numbers
over the field of rational numbers, and
x(1) = (1, 3, 0) = x(1)1 , x(1)2 , x(1)3 (7.6)
1 1
x(2) = , 3, 0 = , 3, 0 (7.7)
1+1 2
1 2
x(3) = 1 , 3, 0 = , 3, 0 (7.8)
1+ 2 3
1 3
x(4) = , 3, 0 = , 3, 0 (7.9)
1 + 23 5
..
. (7.10)
1
x(n) = , 3, 0 (7.11)
1 + x(n−1)1
for n ≥ 2. Does this sequence have a limit point in Q3 ? Is this a Cauchy sequence?
Consider the first term only; the other two are trivial. The series has converged when the nth term
is equal to the (n − 1)th term:
1
x(n−1)1 = .
1 + x(n−1)1
Rearranging, it is found that
x2(n−1)1 + x(n−1)1 − 1 = 0.
Solving, one finds that √
−1 ± 5
x(n−1)1 = .
2
We find from numerical experimentation that it is the “+” root to which x1 converges:
√
5−1
lim x(n−1)1 = .
n→∞ 2
As n → ∞,
√ !
5−1
x(n) → , 3, 0 .
2
Thus, the limit point for this sequence is not is Q3 ; hence the sequence is not convergent. Had the set
been defined in R3 , it would have been convergent.
However, the sequence is a Cauchy sequence. Consider, say ǫ = .01. If we choose, we then find by
numerical experimentation that Nǫ = 4. Choosing, for example m = 5 > Nǫ and n = 21 > Nǫ , we get
5
x(5) = , 3, 0 .
8
10946
x(21) = , 3, 0 .
17711
987
||x(5) − x(21) ||2 = , 0, 0 = 0.00696 < 0.01.
141688 2
This could be generalized for arbitrary ǫ, so the sequence can be shown to be a Cauchy sequence.
Example 7.11
Does the infinite sequence of functions
v = {v1 (t), v2 (t), · · · , vn (t), · · ·} = t(t), t(t2 ), t(t3 ), · · · , t(tn ), · · · ,
converge in L2 [0, 1]? Does the sequence converge in C[0, 1]?
First, check if the sequence is a Cauchy sequence:
s r
Z 1
n+1 m+1 2 1 2 1
lim ||vn (t) − vm (t)||2 = (t −t ) dt = − + = 0.
n,m→∞ 0 2n + 3 m + n + 3 2m + 3
As this norm approaches zero, it will be possible for any ǫ > 0 to find an integer Nǫ such that
||vn (t) − vm (t)||2 < ǫ. So, the sequence is a Cauchy sequence. We also have
0, 0 ≤ t < 1
lim vn (t) =
n→∞ 1, t = 1
The above function, the “limit point” to which the sequence converges, is in L2 [0, 1], which is sufficient
condition for convergence of the sequence of functions in L2 [0, 1]. However the “limit point” is not a
continuous function, so despite the fact that the sequence is a Cauchy sequence and elements of the
sequence are in C[0, 1], the sequence does not converge in C[0, 1].
Example 7.12
Analyze the sequence of functions
n√ √ √ o
v = {v1 , v2 , . . . , vn , . . .} = 2 sin(πt), 2 sin(2πt), . . . , 2 sin(nπt), . . .
in L2 [0, 1].
This is simply a set of sin functions, which can be shown to form a basis; such a proof will not be
given here. Each element of the set is orthonormal to other elements:
Z 1 √ 2 1/2
||vn (t)||2 = 2 sin(nπt) dt = 1.
0
R1
It is also easy to show that 0 vn (t)vm (t) dt = 0, so the basis is orthonormal. As n → ∞, the norm of
the basis function remains bounded, and is, in fact, unity.
Consider the norm of the difference of the mth and nth functions:
Z 1 √ √ 2 21 √
||vn (t) − vm (t)||2 = 2 sin(nπt) − 2 sin(mπt) dt = 2.
0
This is valid for all m and n. Since we can find a value of ǫ > 0 which violates the conditions for a
Cauchy sequence, this series of functions is not a Cauchy sequence.
4. <x, y> = <y, x>, where <·> indicates the complex conjugate of the inner product.
Inner product spaces are subspaces of linear vector spaces and are sometimes called pre-
Hilbert9 spaces. A pre-Hilbert space is not necessarily complete, so it may or may not form
a Banach space.
Example 7.13
Show
<αx, y> = α<x, y>
Using the properties of the inner product and the complex conjugate we have
As a counterexample if x ∈ R2 , and we take ||x|| = ||x||3 = (|x1 |3 + |x2 |3 )1/3 (thus x ∈ ℓ23
which is a Banach space), we cannot find a definition of the inner product which satisfies all
its properties. Thus the space ℓ23 cannot be a Hilbert space!
Unless specified otherwise the unsubscripted norm || · || can be taken to represent the
Hilbert space norm || · ||2 . It is quite common for both subscripted and unscripted versions
of the norm to appear in the literature.
Examples of spaces which are Hilbert spaces are
2. Lebesgue spaces
Rb
• x ∈ L2 [a, b], y ∈ L2 [a, b], t ∈ [a, b] ∈ R1 with <x, y> = a
x(t)y(t) dt.
Rb
• x ∈ L2 [a, b], y ∈ L2 [a, b], t ∈ [a, b] ∈ R1 with <x, y> = a
x(t)y(t) dt.
3. Sobolov spaces
∂u
• u ∈ W12 (G), v ∈ W12 (G), x ∈ G ∈ Rn , n ∈ N, u ∈ L2 (G), ∂x i
∈ L2 (G), v ∈
∂v
L2 (G), ∂xi ∈ L2 (G) with
Z n
!
X ∂u ∂v
<u, v> = u(x)v(x) + dx.
G i=1
∂xi ∂xi
l 2(C 1) complex
scalars
l 2 (C n) n-dimensional
complex vectors
L 2 Lebesgue integrable
function space
1
W
2 Sobolov space
Hilbert space
(complete, normed, inner product)
Banach space
(complete, normed)
Minkowski space
Linear space
Space
Figure 7.3: Venn diagram showing relationship between various classes of spaces.
for complex vectors, we might have taken <x, y> = xT y. Then if we had taken x = (i, 1)T
and y = (1, 1)T , we would have <x, y> = <y, x> = 1 + i. However, we would also have
<x, x> = ||x||22 = (i, 1)(i, 1)T = 0! Obviously, this would violate the property of the norm
since we must have ||x||22 > 0 for x 6= 0.
Interestingly, one can interpret the Heisenberg11 uncertainty principle to be entirely con-
sistent with our definition of an inner product which does not commute in a complex space.
In quantum mechanics, the superposition of physical states of a system is defined by a
complex-valued vector field. Position is determined by application of a position operator,
and momentum is determined by application of a momentum operator. If one wants to know
both position and momentum, both operators are applied. However, they do not commute,
and application of them in different orders leads to a result which varies by a factor related
to Planck’s12 constant.
Matrix multiplicaton is another example of an inner product that does not commute,
in general. Such topics are considered in the more general group theory. Operators that
commute are known as Abelian13 and those that do not are known as non-Abelian.
Obviously, this inner product can take on negative values. The theory goes on to show that
when relativistic effects are important, ordinary concepts of Euclidean geometry become
meaningless, and a variety of non-intuitive results can be obtained. In the Venn diagram,
11
Werner Karl Heisenberg, 1901-1976, German physicist.
12
Max Karl Ernst Ludwig Planck, 1858-1947, German physicist.
13
Niels Henrick Abel, 1802-1829, Norweigen mathematician, considered solution of quintic equations by
elliptic functions, proved impossibility of solving quintic equations with radicals, gave first solution of an
integral equation, famously ignored by Gauss.
we see that Minkowski spaces certainly are not Banach, but there are also linear spaces that
are not Minkowski, so it occupies an island in the diagram.
Example 7.14
For x and y belonging to a Hilbert space, prove the parallelogram equality
<x + y, x + y> + <x − y, x − y> = (<x, x> + <x, y> + <y, x> + <y, y>)
+ (<x, x> − <x, y> − <y, x> + <y, y>)
= 2<x, x> + 2<y, y>
= 2||x||22 + 2||y||22
Example 7.15
Prove the Schwarz14 inequality
Note that this effectively defines the angle between two vectors. Because of the inequality, we have
||x||2 ||y||2
≥ 1,
|<x, y>|
|<x, y>|
≤ 1.
||x||2 ||y||2
Defining α to be the angle between the vectors x and y, we recover the familiar result from vector
analysis
<x, y>
cos α = . (7.18)
||x||2 ||y||2
This reduces to the ordinary relationship we find in Euclidean geometry when x, y ∈ R3 .
Example 7.16
For x, y ∈ ℓ2 (R2 ), find <x, y> if
1 2
x= , y= .
3 −2
The solution is
2
<x, y> = xT y = ( 1 3 ) = (1)(2) + (3)(−2) = −4.
−2
Note that the inner product yields a real scalar, but in contrast
√ √ to the norm, it can be negative. Note
also that the Schwarz inequality holds as ||x||2√||y||2 = 10 8 ∼ 8.944 √ > |√− 4|. Also the Minkowski
inequality holds as ||x + y||2 = ||(3, 1)T ||2 = + 10 < ||x||2 + ||y||2 = 10 + 8.
Example 7.17
For x, y ∈ ℓ2 (C2 ), find <x, y> if
−1 + i 1 − 2i
x= , y= .
3 − 2i −2
The solution is
1 − 2i
<x, y> = xT y = ( −1 − i 3 + 2i ) = (−1 − i)(1 − 2i) + (3 + 2i)(−2) = −9 − 3i.
−2
Note that the inner product is a complex scalar which has negative components. It is easily shown that
||x||2 = 3.870 and ||y||2 = 3 and ||x + y||2 = 2.4495. Also |<x, y>| = 9.4868. The Schwarz inequality
holds as (3.870)(3) = 11.61 > 9.4868. The Minkowski inequality holds as 2.4495 < 3.870 + 3 = 6.870.
Example 7.18
For x, y ∈ L2 [0, 1], find <x, y> if
x(t) = 3t + 4, y(t) = −t − 1.
The solution is
Z 1 1
7t2
3 17
<x, y> = (3t + 4)(−t − 1) dt = −4t − −t =− = −8.5.
0 2 0 2
Once more the inner product is a negative scalar. It is easily shown that ||x||2 = 5.56776 and ||y||2 =
1.52753 and ||x+ y||2 = 4.04145. Also |<x, y>| = 8.5. It is easily seen that the Schwarz inequality holds
as (5.56776)(1.52753) = 8.505 > 8.5. The Minkowski inequality holds as 4.04145 < 5.56776 + 1.52753 =
7.095.
Example 7.19
For x, y ∈ L2 [0, 1], find <x, y> if
x(t) = it, y(t) = t + i.
We recall that Z 1
<x, y> = x(t)y(t) dt.
0
The solution is Z 1
1
t2 it3 1 i
<x, y> = (−it)(t + i) dt = − = − .
0 2 3 0 2 3
The inner product is a complex scalar. It is easily shown that ||x||2 = 0.57735 and ||y||2 = 1.1547
and ||x + y||2 = 1.63299. Also |<x, y>| = 0.601. The Schwarz inequality holds as (0.57735)(1.1547) =
0.6667 > 0.601. The Minkowski inequality holds as 1.63299 < 0.57735 + 1.1547 = 1.7321.
Example 7.20
For u, v ∈ W21 (G)), find <u, v> if
u(x) = x1 + x2 , v(x) = −x1 x2 ,
and G is the square region in the x1 , x2 plane x1 ∈ [0, 1], x2 ∈ [0, 1]. We recall that
Z
∂u ∂v ∂u ∂v
<u, v> = u(x)v(x) + + dx,
G ∂x1 ∂x1 ∂x2 ∂x2
Z 1Z 1
4
<u, v> = ((x1 + x2 )(−x1 x2 ) + (1)(−x2 ) + (1)(−x1 )) dx1 dx2 = − = −1.33333.
0 0 3
The inner product here is negative real scalar. It is easily shown that ||u||1,2 = 1.77951 and ||v||1,2 =
0.881917 and ||u + v||1,2 = 1.13039. Also |<u, v>| = 1.33333. The Schwarz inequality holds as
(1.77951)(0.881917) = 1.56938 > 1.33333. The Minkowski inequality holds as 1.13039 < 1.77951 +
0.881917 = 2.66143.
7.3.2.4 Orthogonality
One of the primary advantages of working in Hilbert spaces is that the inner product allows
one to utilize of the useful concept of orthogonality:
• x and y are said to be orthogonal to each other if
<x, y> = 0
• In an orthogonal set of vectors {v1 , v2 , · · ·} the elements of the set are all orthogonal
to each other, so that <vi , vj > = 0 if i 6= j.
• If a set {ϕ1 , ϕ2 , · · ·} exists such that <ϕi , ϕj > = δij , then the elements of the set are
orthonormal.
• A basis {v1 , v2 , · · · , vn } of a finite-dimensional space that is also orthogonal is an or-
thogonal basis. On dividing each vector by its norm we get
vi
ϕi = √
<vi , vi >
to give us an orthonormal basis {ϕ1 , ϕ2 , · · · , ϕn }.
Example 7.21
If elements x and y of an inner product space are orthogonal to each other, prove the Pythagorean
theorem
||x||22 + ||y||22 = ||x + y||22
The right side is
<x + y, x + y> = <x, x> + <x, y> + <y, x> + <y, y>
= <x, x> + <y, y> since <x, y> = <y, x> = 0 due to orthogonality
= ||x||22 + ||y||22 QED
Example 7.22
Show that an orthogonal set of vectors in an inner product space is linearly independent.
Let {v1 , v2 , · · ·} be an orthogonal set of vectors. Then consider
α1 v1 + α2 v2 + . . . + αj vj + . . . + αn vn = 0.
Taking the inner product with vj , where j = 1, 2, . . . we get
<vj , (α1 v1 + α2 v2 + . . . + αj vj + . . . + αn vn )> = <vj , 0>
α1 <vj , v1 > + α2 <vj , v2 > + . . . + αj <vj , vj > + . . . + αn <vj , vn > = 0
αj <vj , vj > = 0
since all the other inner products are zero. Thus, αj = 0, indicating that the set {v1 , v2 , · · ·} is linearly
independent.
Example 7.23
Find an orthonormal set of vectors {ϕ1 , ϕ2 , . . .} in L2 [−1, 1] using linear combinations of the linearly
independent set of vectors {1, t, t2 , t3 , . . .} where −1 ≤ t ≤ 1.
Choose
v1 (t) = 1
Now choose the second vector linearly independent of v1 as
v2 (t) = a + bt.
This should be orthogonal to v1 , so that
Z 1
v1 (t)v2 (t) dt = 0
−1
Z 1
(1) (a + bt) dt = 0
−1 |{z} | {z }
=v1 (t) =v2 (t)
1
bt2
at + =0
2 −1
b
a(1 − (−1)) + (12 − (−1)2 ) = 0
2
from which
a=0
Taking b = 1 arbitrarily, since orthogonality does not depend on the magnitude of v2 (t), we have
v2 = t.
Choose the third vector linearly independent of v1 (t) and v2 (t), i.e.
v3 (t) = a + bt + ct2 .
For this to be orthogonal to v1 (t) and v2 (t), we get the conditions
Z 1
(1) (a + bt + ct2 ) dt = 0
−1 |{z} | {z }
=v1 (t) =v3 (t)
Z 1
t (a + bt + ct2 ) dt
|{z} = 0
−1 | {z }
=v2 (t) =v3 (t)
15
Jørgen Pedersen Gram, 1850-1916, Danish mathematician, and Erhard Schmidt, 1876-1959,
German/Estonian-born Berlin mathematician, studied under David Hilbert, founder of modern functional
analysis. The Gram-Schmidt procedure was actually first introduced by Laplace.
The first of these gives c = −3a. Taking a = 1 arbitrarily, we have c = −3. The second relation gives
b = 0. Thus
v3 = 1 − 3t2
In this manner we can find as many orthogonal vectors as we want. We can make them orthonormal
by dividing each by its norm, so that we have
1
ϕ1 = √
2
r
3
ϕ2 = t
2
r
5
ϕ3 = (1 − 3t2 )
8
..
.
Scalar multiples of these functions, with the functions set to unity at t = 1, are the Legendre polyno-
mials: P0 (t) = 1, P1 (t) = t, P2 (t) = 12 (3t2 − 1) . . .
2.
ϕn (t) = e−t/2 Ln (t),
where Ln (t) are the Laguerre polynomials for t ∈ [0, ∞).
x = a1 u1 + a2 u2 + · · · + an un (7.19)
The general task here is to find expressions for the coefficients ak , k = 1, 2, . . . n. To get the
coefficients, we begin by taking inner products with u1 , u2 , · · · , un in turn to get
Knowing x, u1 , u2, · · · , un , all the inner products can be determined, and the equations can
be posed as a linear algebraic system:
<u1 , u1> <u1 , u2 > . . . <u1 , un > a1 <u1 , x>
<u2 , u1> <u2 , u2 > . . . <u2 , un > a2 <u2 , x>
.. .. .. . = .. (7.20)
. . ... . .. .
<un , u1 > <un , u2 > . . . <un , un > an <un , x>
This can also be written as
<ui, uj >aj = <ui, x>
In either case, Cramer’s rule can be used to solve for the unknown coefficients, aj .
The process is simpler if the basis vectors are orthogonal. If orthogonal,
<ui , uj > = 0, i 6= j
If we use an orthonormal basis {ϕ1 , ϕ2 , . . . , ϕn } then the representation is even more efficient:
n
X
x= <ϕi , x> ϕi (7.25)
i=1
Similar expansions apply to vectors in infinite-dimensional spaces, except that one must
be careful that the orthonormal set is complete. Only then is there any guarantee that any
vector can be represented as linear combinations of this orthonormal set.
If {ϕ1 , ϕ2 , . . .} is a complete orthonormal set of vectors in some domain Ω, then any
vector x can be represented as
X∞
x= ai ϕi (7.26)
i=1
where
ai = <ϕi , x> (7.27)
This is a Fourier series representation, and the ai are the Fourier coefficients. Though
trigonometric functions are sometimes used, other orthogonal functions are also common.
Thus we can have Fourier-Legendre for Ω = [−1, 1], Fourier-Hermite for Ω = (−∞, ∞), or
Fourier-Laguerre series for Ω = [0, ∞).
Example 7.24
Expand the top hat function x(t) = H(t − 1/4) − H(t − 3/4) in a Fourier sine series in the domain
t ∈ [0, 1].
Here, the function x(t) is discontinuous at t = 1/4 and t = 3/4. While x(t) is not a member of
C[0, 1], it is a member of L2 [0, 1]. Here we will see that the Fourier sine series representation, composed
of functions which are continuous in [0, 1], converges to the discontinuous function x(t).
Building on previous work, we know that the set of functions
√
ϕn (t) = 2 sin(nπt), n = 1, . . . , ∞,
are orthonormal for t ∈ [0, 1]. We then find for the Fourier coefficients
√ Z 1
1 3 √ Z 3/4
an = 2 H t− −H t− sin(nπt) dt = 2 sin(nπt) dt.
0 4 4 1/4
The discontinuous function x(t) and three continuous approximations to it are shown in Figure 7.4.
Note that as more terms are added, the approximation gets better at most points. But there is always
a persistently large error at the discontinuities t = 1/4, t = 3/4. We say this function is convergent in
L2 [0, 1], but is not convergent in L∞ [0, 1]. This simply says that the rms error norm converges, while the
x (t)
1
9 term series
0.8
0.6
0.4
0.2
t
0.2 0.4 0.6 0.8 1
x (t)
1
18 term series
0.8
0.6
0.4
0.2
t
0.2 0.4 0.6 0.8 1
x (t)
1
36 term series
0.8
0.6
0.4
0.2
t
0.2 0.4 0.6 0.8 1
Figure 7.4: Expansion of top hat function x(t) = H(t − 1/4) − H(t − 3/4) in terms of sine
basis functions for three levels of approximation, n = 9, n = 18, n = 36.
|| x a(t) - x (t) ||
2
0.70
-0.512
0.50 || x a(t) - x (t) || ~ 0.474 n
2
0.30
0.20
0.15
0.10
2 5 10 20 n
maximum error norm does not. This is an example of the well-known Gibbs phenomenon. Convergence
in L2 [0, 1] is shown in Fig. 7.5. The achieved convergence rate is ||x(t) − xa (t)||2 ∼ 0.474088n−0.512.
This suggests that
1
lim ||x(t) − xa (t)||2 ∼ √ .
n→∞ n
The previous example showed one could use continuous functions to approximate a dis-
continuous function. The converse is also true: discontinuous functions can be used to
approximate continuous functions.
Example 7.25
Show that the functions ϕ1 (t), ϕ2 (t), . . . , ϕn (t) are orthonormal in L2 (0, 1], where
√
n k−1 n < t≤ n
k
ϕk (t) =
0 otherwise
Expand x(t) = t2 in terms of these functions, and find the error for a finite n.
We note that the basis functions are a set of “top hat” functions whose amplitude increases with
and whose width decreases with increasing n. For fixed n, the basis functions √ are a series of top hats
that fills the domain [0, 1]. The area enclosed by a single basis function is 1/ n. If k 6= j the inner
product
Z 1
<ϕk , ϕj > = ϕk (t)ϕj (t) dt = 0
0
So, ϕ1 , ϕ2 , . . . , ϕn is an orthonormal set. We can expand the function f (t) = t2 in the form
n
X
t2 = αk ϕk
k=1
x(t) x(t) 2
2
x(t) = t x(t) = t
1 1
n=5
0.8 0.8 n = 10
0.6 0.6
0.4 0.4
0.2 0.2
t t
0.2 0.4 0.6 0.8 1 0.2 0.4 0.6 0.8 1
Figure 7.6: Expansion of x(t) = t2 in terms of “top hat” basis functions for two levels of
approximation, n = 5, n = 10.
Z 1 n
X
ϕj (t)t2 dt = αk δkj
0 k=1
Z 1
ϕj (t)t2 dt = αj
0
Z 1
ϕk (t)t2 dt = αk
0
Thus Z k
n √
αk = 0 + t2 n dt + 0
k−1
n
Thus
1
αk = 5/2
3k 2 − 3k − 1
3n
2
Pn
The functions t and the partial sums fn (t) = k=1 αk ϕk (t) for n = 5 and n = 10 are shown in
Figure 7.6. The L2 error for the partial sums can be calculated as ∆n , where
∆2n = ||f (t) − fn (t)||22
Z 1 n
!2
X
2
= t − αk ϕk (t) dt
0 k=1
1 1
= 1− 2
9n2 5n
r
1 1
∆n = 1− 2
3n 5n
1
which vanishes as n → ∞ at a rate of convergence proportional to n.
Example 7.26
Show the Fourier sine series for x(t) = 2t converges at a rate proportional to √1 , where n is the
n
number of terms used to approximate x(t), in L2 [0, 1].
It is easy to show linear independence for these functions. They are orthonormal in the Hilbert space
L2 [0, 1], e.g.
Z 1 √ √
<ϕ2 , ϕ3 > = 2 sin(2πt) 2 sin(3πt) dt = 0,
0
Z 1 √ √
<ϕ3 , ϕ3 > = 2 sin(3πt) 2 sin(3πt) dt = 1.
0
Note that while the basis functions evaluate to 0 at both t = 0 and t = 1, that the function itself
only has value 0 at t = 0. We must tolerate a large error at t = 1, but hope that this error is confined
to an ever collapsing neighborhood around t = 1 as more terms are included in the approximation.
The Fourier coefficients are
Z 1√ √
2 2(−1)k+1
αk = <ϕk (t), 2t> = 2 sin(kπt) (2t) dt = .
0 kπ
The approximation then is
n
X 4(−1)k+1
xa (t) = sin(kπt).
kπ
k=1
This is difficult to evaluate analytically. It is straightforward to examine this with symbolic calculational
software.
A plot of the norm of the error as a function of the number of terms in the approximation, n,
is given in the log-log plot of Figure 7.7. A weighted least squares curve fit, with a weighting factor
proportional to n2 so that priority is given to data as n → ∞, shows that the function
approximates the convergence performance well. In the log-log plot the exponent on n is the slope. It
appears from the graph that the slope may be approaching a limit, in which it is likely that
1
||x(t) − xa (t)||2 ∼ √ .
n
This indicates convergence of this series. Note that the series converges even though the norm of the
nth basis function does not approach zero as n → ∞:
since the basis functions are orthonormal. Also note that the behavior of the norm of the final term in
the series, v
u 1 2√2(−1)n+1 √ √
uZ !2
t 2 2
||αn ϕn (t)||2 = 2 sin(nπt) dt = ,
0 nπ nπ
does not tell us how the series actually converges.
|| x(t) - x a(t) ||
2
0.7
- 0.481
|| x(t) - x a(t) || ~ 0.841 n
2
0.5
0.3
0.2
1 1.5 2 3 5 7 10 15 20 n
Figure 7.7: Behavior in the error norm of the Fourier series approximation to x(t) = 2t with
the number n of terms included in the series.
Example 7.27
1
Show the Fourier sine series for x(t) = t − t2 converges at a rate proportional to n5/2
, where n is
the number of terms used to approximate x(t), in L2 [0, 1].
Again, consider the sequence of functions
n√ √ √ o
ϕ= 2 sin(πt), 2 sin(2πt), . . . , 2 sin(nπt), . . . .
which are as before, linearly independent and moreover, orthonormal. Note that in this case, as opposed
to the previous example, both the basis functions and the function to be approximated vanish identically
at both t = 0 and t = 1. Consequently, there will be no error in the approximation at either end point.
The Fourier coefficients are
√
2 2 1 + (−1)k+1
αk = .
k3 π3
Note that αk = 0 for even values of k. Taking this into account and retaining only the necessary basis
functions, we can write the Fourier sine series as
n √
X 4 2
x(t) = t(1 − t) ∼ xa (t) = sin((2m − 1)πt).
m=1
(2m − 1)3 π 3
Again this is difficult to address analytically, but symbolic computation allows computation of the error
norm as a function of n.
|| x(t) - xa(t)||
2
0.005
0.001 - 2.492
0.0005 || x(t) - xa(t)|| ~ 0.00994 n
2
0.0001
0.00005
0.00001
1 1.5 2 3 5 7 10 15 20 n
Figure 7.8: Behavior in the error norm of the Fourier series approximation to x(t) = t(1 − t)
with the number n of terms included in the series.
A plot of the norm of the error as a function of the number of terms in the approximation, n,
is given in the log-log plot of Figure 7.8. A weighted least squares curve fit, with a weighting factor
proportional to n2 so that priority is given to data as n → ∞, shows that the function
||x(t) − xa (t)||2 ∼ 0.00995 n−2.492 ,
approximates the convergence performance well. Thus we might suspect that
1
lim ||x(t) − xa (t)||2 ∼ .
n→∞ n5/2
Note that the convergence is much more rapid than in the previous example! This can be critically
important in numerical calculations and demonstrates that a judicious selection of basis functions can
have fruitful consequences.
Here we have used the definition of the Fourier coefficient <ϕ, x> = α, and orthonormality <ϕ, ϕ> = 1.
This is easily extended to multiterm expansions to give
2
n
X n
X
2
|αk |2 .
x(t) − α ϕ
k k (t) = ||x(t)||2 − (7.35)
k=1 2 k=1
for all functions x(t). Note that this requirement is stronger than just requiring that the last Fourier
coefficient vanish for large n; also note that it does not address the important question of the rate of
convergence, which can be different for different functions x(t), for the same basis.
<ui , uR
j > = δij
Then {uR R
1 , · · · , un } is called the reciprocal (or dual) basis of {u1 , · · · , un }. Of course an orthonormal
basis is its own reciprocal.
Since {u1 , · · · , un } is a basis, we can write any vector x as
n
X
x= αj uj (7.37)
j=1
n
X
<uR
i , x> = <uR
i , αj uj > (7.38)
j=1
n
X
= <uR
i , αj uj > (7.39)
j=1
Xn
= αj <uR
i , uj > (7.40)
j=1
Xn
= αj δij (7.41)
j=1
= αi (7.42)
<uR
j , x> = αj (7.43)
so that
n
X
x= <uR
j , x> uj . (7.44)
j=1
| {z }
=αj
Example 7.28
2 1
Consider x ∈ R2 . The vectors u1 = and u2 = span the space R2 and thus can be used
0 3
as a basis.
3
Find the reciprocal basis uR R
1 , u2 , and use the above relation to expand x = in terms of both
5
the basis u1 and u2 and then the reciprocal basis uR R
1 and u2 .
We adopt the dot product as our inner product. Let’s get α1 , α2 . To do this we first need the
reciprocal basis vectors which are defined by the inner product:
<ui , uR
j > = δij
We take
a R a2R 1
uR
1 =
1 1
, uR2 =
a1R 2 a2R 2
expanding, we get
R T R a1R 1
<u1 , u1 > = u1 u1 = (2, 0) · = (2)a1R 1 + (0)a1R 2 = 1
a1R 2
R T R a2R 1
<u1 , u2 > = u1 u2 = (2, 0) · = (2)a2R 1 + (0)a2R 2 = 0
a2R 2
R T R a1R 1
<u2 , u1 > = u2 u1 = (1, 3) · = (1)a1R 1 + (3)a1R 2 = 0
a1R 2
R T R a2R 1
<u2 , u2 > = u2 u2 = (1, 3) · = (1)a2R 1 + (3)a2R 2 = 1
a2R 2
Solving, we get
1 1 1
a1R 1 = , a1R 2 = − , a2R 1 = 0, a2R 2 =
2 6 3
so substituting, we get expressions for the reciprocal base vectors:
1
0
uR = 2 , u R
=
1 − 16 2 1
3
x
2
2/3 u 1
u2
R
18 u2
5/3 u 2
u2R
u1
uR
1
x1
6 uR
1
In a similar manner it is easily shown that x can be represented in terms of the reciprocal basis as
n
X
x= βi u R R R
i = β1 u 1 + β2 u 2 ,
i=1
where
βi = <ui , x>.
For this problem, this yields
x = 6uR R
1 + 18u2 .
Thus we see for the non-orthogonal basis that two natural representations of the same vector exist.
One of these is actually a a covariant representation; the other is contravariant.
Let us show this example is consistent with the earlier described notions using “upstairs-downstairs”
index notation. Note that our non-orthogonal coordinate system is a transformation of the form
∂ξ i j
ξi = x ,
∂xj
where ξ i is the Cartesian representation, and xj is the contravariant representation in the transformed
system. In Gibbs form, this is
ξ = J · x.
so that
1
ξ1 2 1 x
= · .
ξ2 0 3 x2
has representation in Cartesian space of (2, 0)T , and the other unit vector in the transformed space
x1 0
=
x2 1
The Cartesian vector ξ = (3, 5)T , has a contravariant representation in the transformed space of
−1
2 1 3 1/2 −1/6 3 2/3
x = J−1 · ξ = · = · = = xj .
0 3 5 0 1/3 5 5/3
The rows of this matrix describe the reciprocal basis vectors, and is also consistent with our earlier
finding.
7.4 Operators
• For two sets X and Y, an operator (or mapping, or transformation) f is a rule that
f
associates every x ∈ X with an image y ∈ Y. We can write f : X → Y, X → Y or
x 7→ y. X is the domain of the operator, and Y is the range.
• If every element of Y is not necessarily an image, then X is mapped into Y; this map
is called an injection.
• If, on the other hand, every element of Y is an image of some element of X, then X is
mapped onto Y and the map is a surjection.
• If, for every x ∈ X there is a unique y ∈ Y, and for every y ∈ Y there is a unique
x ∈ X, the operator is one-to-one or invertible; it is a bijection.
f g
• f and g are inverses of each other, when X → Y and Y → X.
T(x + y) = Tx + Ty (7.45)
T(αx) = αTx (7.46)
. ..
Domain Injection: Inverse may not exist Range
. f
f
. . Y
X
.. . ..
Surjection: Inverse not always unique
f
f
X f Y
... f
f
.. .
Bijection (one-to-one): Inverse always exists
X f Y
• The null space or kernel of an operator T is the set of all x such that Tx = 0. The
null space is a vector space.
• An operator T is
∀ x 6= 0.
• Theorem
√
For a matrix A, Cm → Cn , ||A||2 = λmax , where λmax is the largest eigenvalue of
T T
the matrix A A. It will soon be shown that because A A is symmetric, that all of
its eigenvalues are guaranteed real. Moreover, it can be shown that they are also all
greater than or equal to zero. Hence, the definition will satisfy all properties of the
norm.
• the above theorem holds only for Hilbert spaces and not for arbitrary Banach spaces.
Example 7.29
Find the adjoint of the real matrix A : R2 → R2 , where
a11 a12
A=
a21 a22
Example 7.30
Find the adjoint of the differential operator L : X → X, where
d2 d
L= 2
+
ds ds
This maintains the form of an inner product in L2 [0, 1] if we require y(0) = y(1) = 0; doing this, we get
Z 1
<Lx, y> = x(s) (y ′′ (s) − y ′ (s)) ds = <x, L∗ y>
0
Example 7.31
d2
Find the adjoint of the differential operator L : X → X, where L = ds2 , and X is the subspace of
L2 [0, 1] with x(0) = x(1) = 0 if x ∈ X.
Using integration by parts on the inner product
Z 1
<Lx, y> = x′′ (s)y(s) ds
0
Z 1
′ ′
= x (1)y(1) − x (0)y(0) − x′ (s)y ′ (s) ds
0
Z 1
= x′ (1)y(1) − x′ (0)y(0) − x(1) y ′ (1) − x(0) y ′ (0) − x(s)y ′′ (s) ds
|{z} |{z} 0
=0 =0
Z 1
= x′ (1)y(1) − x′ (0)y(0) + x(s)y ′′ (s) ds
0
Example 7.32
Find the adjoint of the integral operator L : L2 [a, b] → L2 [a, b], where
Z b
Lx = K(s, t)x(s) ds.
a
The inner product
Z Z !
b b
<Lx, y> = K(s, t)x(s) ds y(t) dt
a a
Z b Z b
= K(s, t)x(s)y(t) ds dt
a a
Z b Z b
= x(s)K(s, t)y(t) dt ds
a a
Z Z !
b b
= x(s) K(s, t)y(t) dt ds
a a
= <x, L∗ y>
where Z b
∗
L y= K(s, t)y(t) dt
a
or equivalently
Z b
L∗ y = K(t, s)y(s) ds
a
Note in the definition of Lx, the second argument of K is a free variable, while in the consequent
definition of L∗ y, the first argument of K is a free argument. So in general, the operator and its adjoint
are different. Note however, that
if K(s, t) = K(t, s) then the operator is self-adjoint
That is, a symmetric kernel yields a self-adjoint operator.
Properties:
||T∗|| = ||T|| (7.50)
(T1 + T2 )∗ = T∗1 + T∗2 (7.51)
(αT)∗ = αT∗ (7.52)
(T1 T2 )∗ = T∗2 T∗1 (7.53)
(T∗ )∗ = T (7.54)
(T−1 )∗ = (T∗ )−1 if T−1 exists (7.55)
x = T−1 y
TT−1 y = y
so that
TT−1 = I
Properties:
T−1
a y = Tb x, (7.58)
−1 −1
Tb Ta y = x. (7.59)
Example 7.33
Let L be the operator defined by
d2 2
Lx = + k x(t) = f (t)
dt2
where x belongs to the subspace of L2 [0, π] with x(0) = a and x(π) = b. Show that the inverse operator
L−1 is given by
Z π
−1 ∂g ∂g
x(t) = L f (t) = b (π, t) − a (0, t) + g(τ, t)f (τ ) dτ
∂τ ∂τ 0
Using integration by parts and the property that g(0, t) = g(π, t) = 0, the integral in the right can be
simplified as Z π 2
∂g ∂g ∂ g 2
− x(π) (π, t) + x(0) (0, t) + x(τ ) + k g dτ
|{z} ∂τ |{z} ∂τ 0 ∂τ 2
=b =a | {z }
=δ(t−τ )
Thus, L−1 L = I, proving the proposition. Note, it is easily shown for this problem that the Green’s
function is
sin(k(π − τ )) sin(kt)
g(τ, t) = − t<τ
k sin(kπ)
sin(kτ ) sin(k(π − t))
= − τ <t
k sin(kπ)
so that we can write x(t) explicitly in terms of the forcing function f (t) including the inhomogeneous
boundary conditions as follows:
For linear algebraic systems, the reciprocal or dual basis can be easily formulated in
terms of operator notation and is closely related to the inverse operator. If we define U
to be a n × n matrix which has the n basis vectors ui, each of length n, which span the
n-dimensional space, we seek UR , the n × n matrix which has as its columns the vectors
uR
j which form the reciprocal or dual basis. The reciprocal basis is found by enforcing the
equivalent of <ui , uR
j > = δij :
T
U · UR = I. (7.61)
Solving for UR ,
T
U · UR = I, (7.62)
R
UT · U = I, (7.63)
R T
T
U ·U = IT , (7.64)
RT
U · U = I, (7.65)
RT
U · U · U−1 = I · U−1 , (7.66)
RT
U = U−1 , (7.67)
T
UR = U−1 , (7.68)
we see that the set of reciprocal basis vectors is given by the conjugate transpose of the inverse
of the original matrix of basis vectors. Then the expression for the amplitudes modulating
the basis vectors, αi = <uR i , x>, is
T
α = UR · x. (7.69)
Substituting for UR in terms of its definition, we can also say
T
T
· x = U−1 · x.
α = U−1 (7.70)
Pn Pn R
Then the expansion for the vector x = j=1 αj uj = j=1 <uj , x>uj is written in the
alternate notation as
x = U · α = U · U−1 · x = x. (7.71)
Example 7.34
2
Consider the problem of a previous example with x ∈ R2 and with basis vectors u1 = and
0
1
u2 = , find the reciprocal basis vectors and an expansion of x in terms of the basis vectors.
3
Using the alternate vector and matrix notation, we define the matrix of basis vectors as
2 1
U= .
0 3
Since this matrix is real, the complex conjugation process is not important, but it will be retained for
completeness. Using standard techniques, we find that the inverse is
1
− 16
U−1 = 2 1 .
0 3
Thus the matrix with the reciprocal basis vectors in its columns is
1
R −1
T
2 0
U =U = .
− 61 13
This agrees with the earlier analysis. For x = (3, 5)T , we find the coefficients α to be
1 2
R
T
2 − 61 3
α=U ·x= 1 · = 35 .
0 3 5 3
We see that we do indeed recover x upon taking the product
2
2 1 2 2 5 1 3
x=U·α= · 35 = + = .
0 3 3 3 0 3 3 5
Theorem
The eigenvalues of a self-adjoint operator are real.
Proof:
Since the operator is self-adjoint, we have
Here we note that for non-trivial eigenvectors <e, e> > 0, so the division can be performed.
The only way a complex number can equal its conjugate is if its imaginary part is zero;
consequently, the eigenvalue must be strictly real.
Theorem
The eigenvectors of a self-adjoint operator corresponding to distinct eigenvalues are or-
thogonal.
since λi 6= λj .
Theorem
The eigenvectors of any self-adjoint operator on vectors of a finite-dimensional vector
space constitute a basis for the space.
As discussed by Friedman, the following conditions are sufficient for the eigenvectors in
a infinite-dimensional Hilbert space to be form a complete basis:
If the operator is not self-adjoint, Friedman (p. 204) discusses how the eigenfunctions of
the adjoint operator can be used to obtain the coefficients αk on the eigenfunctions of the
operator.
Example 7.35
For x ∈ R2 , A : R2 → R2 , Find the eigenvalues and eigenvectors of
2 1
A=
1 2
then
2−λ 1 x1 0
= (a)
1 2−λ x2 0
By Cramer’s rule we could say
0 1
det
0 2−λ 0
x1 = = ,
2−λ 1 2−λ 1
det det
1 2−λ 1 2−λ
2−λ 0
det
1 0 0
x2 = = .
2−λ 1 2−λ 1
det det
1 2−λ 1 2−λ
An obvious, but uninteresting solution is the trivial solution x1 = 0, x2 = 0. Nontrivial solutions of x1
and x2 can be obtained only if
2−λ 1
=0
1 2−λ
which gives the characteristic equation
(2 − λ)2 − 1 = 0
Solutions are λ1 = 1 and λ2 = 3. The eigenvector corresponding to each eigenvalue is found in the
following manner. The eigenvalue is substituted in equation (a). A dependent set of equations in x1
and x2 is obtained. The eigenvector solution is thus not unique.
For λ = 1, equation (a) gives
2−1 1 x1 1 1 x1 0
= = ,
1 2−1 x2 1 1 x2 0
−x1 + x2 = 0
1
e2 = β
1
We take β = 1, so that
1
e2 = .
1
Comments:
1. Since the matrix is symmetric (thus self-adjoint), the eigenvalues are real, and the eigenvectors are
orthogonal.
2. We have actually solved for the right eigenvectors.
This is the usual set of eigenvectors. The left eigenvectors can be found from xT A = xT Iλ. Since here
A is equal to its conjugate transpose, xT A = Ax, so the left eigenvectors are the same as the right
eigenvectors. More generally, we can say the left eigenvectors of an operator are the right eigenvectors
T
of the adjoint of that operator, A .
3. Multiplication of an eigenvector by any scalar is also an eigenvector.
4. The normalized eigenvectors are
! !
√1 √1
2 2
e1 = , e2 =
− √12 √1
2
Example 7.36
For x ∈ C2 , A : C2 → C2 , find the eigenvalues and eigenvectors of
0 −2
A=
2 0
(A − λI) e = 0.
λ2 + 4 = 0,
which has two imaganary roots which are complex conjugates: λ1 = 2i, λ2 = −2i. The corresponding
eigenvectors are
i −i
e1 = α , e2 = β ,
1 1
Note that
T −i
<e1 , e2 > = e1 e2 = ( −i 1 ) = (−1) + 1 = 0,
1
so this is an orthogonal system. We can render it orthonormal by scaling by the magnitude of each
eigenvector. The orthonormal eigenvector set is
! !
√i − √i2
2
e1 = 1 , e2 = 1 .
√ √
2 2
One should be able to prove that the eigenvectors of an arbitrary anti-symmetric matrix are orthogonal.
Example 7.37
For x ∈ C2 , A : C2 → C2 , find the eigenvalues and eigenvectors of
1 −1
A=
0 1
(A − λI) e = 0.
(1 − λ)2 = 0
which has repeated roots λ = 1, λ = 1. For this eigenvalue, the components of the eigenvector satisfy
the equation
x2 = 0
Thus only one ordinary eigenvector
1
e1 = α
0
can be found. We take arbitrarily α = 1 so that
1
e1 = .
0
(A − λI)g1 = e1 .
(A − λI)2 g1 = 0.
Now
0 −1
(A − λI) =
0 0
So with g1 = (β, γ)T , take
0 −1 β 1
= |{z}
1 .
0 0 γ 0
| {z } | {z } =λ | {z }
=A−λI =g1 =e1
1
We get a solution if β ∈ R , γ = −1. That is
β
g1 = .
−1
Note that the ordinary eigenvector and the generalized eigenvector combine to form a basis, in this case
an orthonormal basis.
Example 7.38
For x ∈ C2 , A : C2 → C2 , find the eigenvalues, right eigenvectors, and left eigenvectors if
1 2
A= .
−3 1
AeR = λIeR .
Note as the operator is not self-adjoint, we are not guaranteed real eigenvalues. The right eigenvectors
are not orthogonal as e1R T e2R = 31 .
For the left eigenvectors, we have
eTL A = eTL Iλ.
We can put this in a slightly more standard form by taking the conjugate transpose of both sides:
T T
eTL A = eTL Iλ .
T
A eL = IλeL .
T
A eL = IλeL .
A∗ eL = Iλ∗ eL .
So the left eigenvectors of A are the right eigenvectors of the adjoint of A. Now we have
T 1 −3
A = .
2 1
Note that in addition to being complex conjugates of themselves, which does not hold for general
complex matrices, the eigenvalues of the adjoint are complex conjugates of those of the original matrix,
which does hold for general complex matrices. That is λ∗ = λ. The left eigenvectors are not orthogonal
as e1L T e2L = − 21 . It is easily shown by taking the conjugate transpose of the adjoint eigenvalue problem
however that
eTL A = eTL λ,
as desired. Note that the eigenvalues for both the left and right eigensystems are the same.
Example 7.39
Consider a small change from the previous example. For x ∈ C2 , A : C2 → C2 , find the eigenvalues,
right eigenvectors, and left eigenvectors if
1 2
A= .
−3 1 + i
AeR = λeR .
λ2 − (2 + i)λ + (7 + i) = 0,
Note as the operator is not self-adjoint, we are not guaranteed real eigenvalues. The right eigenvectors
are not orthogonal as e1R T e2R = 1 6= 0
For the left eigenvectors, we solve the corresponding right eigensystem for the adjoint of A which
T
is A∗ = A .
T 1 −3
A = .
2 1−i
T
The eigenvalue problem is A eL = λ∗ eL . The eigensystem is
∗ 3i ∗ −i
λ1 = 1 + 2i, e1L = ; λ2 = 1 − 3i, e2L = .
2 1
Note that here, the eigenvalues λ∗1 , λ∗2 have no relation to each other, but they are complex conjugates
of the eigenvalues, λ1 , λ2 , of the right eigenvalue problem of the original matrix. The left eigenvectors
are not orthogonal as e1L T e2L = −1. It is easily shown however that
eTL A = eTL λ,
as desired.
Example 7.40
For x ∈ R3 , A : R3 → R3 , find the eigenvalues and eigenvectors of
2 0 0
A= 0 1 1
0 1 1
From
2−λ 0 0
0
1 − λ 1 =0
0 1 1−λ
the characteristic equation is
(2 − λ) (1 − λ)2 − 1 = 0
The solutions are λ = 0, 2, 2. The second eigenvalue is of multiplicity two. Next we find the eigenvectors
x1
e = x2
x3
from which
0
e1 = α 1
−1
For λ = 2, we have
0 0 0 x1 0
0 −1 1 x2 = 0
0 1 −1 x3 0
satisfies the equation. Here, we have two free parameters, β and γ; we can thus extract two independent
eigenvectors from this. For e2 we arbitrarily take β = 0 and γ = 1 to get
0
e2 = 1 .
1
In this case e1 , e2 , e3 are orthogonal even though e2 and e3 correspond to the same eigenvalue.
Example 7.41
For y ∈ L2 [0, 1], find the eigenvalues and eigenvectors of L = −d2 /dt2 , operating on functions which
vanish at 0 and 1.
The eigenvalue problem is
d2 y
Ly = − = λy with y(0) = y(1) = 0
dt2
or
d2 y
+ λy = 0 with y(0) = y(1) = 0
dt2
The solution of this differential equation is
The boundary condition y(0) = 0 gives b = 0. The other condition y(1) = 0 gives a sin λ1/2 = 0. A
nontrivial solution can only be obtained if
sin λ1/2 = 0
There are an infinite but countable number of values of λ for which this can be satisfied. These are
λn = n2 π 2 , n = 1, 2, · · ·. The eigenvectors (also called eigenfunctions in this case) yn (t), n = 1, 2, · · ·
are
yn (t) = sin nπt
The differential operator is self-adjoint so that the eigenvalues are real and the eigenfunctions are
orthogonal.
Example 7.42
For x ∈ L2 [0, 1], and L = d2 /ds2 + d/ds with x(0) = x(1) = 0, find the Fourier expansion of an
arbitrary function f (s) in terms of the eigenfunctions of L. Find the series representation of the “top
hat” function
1 3
f (s) = H s − −H s− .
4 4
We seek expressions for αn in
N
X
f (s) = αn xn (s).
n=1
d2 x dx
Lx = + = λx; x(0) = x(1) = 0.
ds2 ds
It is easily shown that the eigenvalues of L are given by
1
λn = − − n2 π 2 , n = 1, 2, 3, . . .
4
where n is a positive integer, and the unnormalized eigenfunctions of L are
Although the eigenvalues are real, the eigenfunctions are not orthogonal. We see this, for example,
by forming <x1 , x2 >:
Z 1
<x1 , x2 > = e−s/2 sin (πs) e−s/2 sin (2πs) ds,
0 | {z }| {z }
=x1 (s) =x2 (s)
4(1 + e)π 2
<x1 , x2 > = 6= 0.
e(1 + π 2 )(1 + 9π 2 )
By using integration by parts, we calculate the adjoint operator to be
d2 y dy
L∗ y = − = λ∗ y; y(0) = y(1) = 0.
ds2 ds
We then find the eigenvalues of the adjoint operator to be the same as those of the operator (this is
true because the eigenvalues are real; in general they are complex conjugates of one another).
1
λ∗m = λm = − − m2 π 2 , m = 1, 2, 3, . . .
4
where m is a positive integer.
The unnormalized eigenfunctions of the adjoint are
So, for m = n, we get <yn , xn > 6= 0, and for m 6= n, we get <ym , xn > = 0. Then, we must have the
so-called bi-orthogonality condition
Dmn = 0 if m 6= n.
Here Dmn diagonal matrix which can be reduced to the identity matrix with proper normalization.
Now consider the following series of operations on the original form of the expansion we seek
N
X
f (s) = αn xn (s),
n=1
N
X
<yj (s), f (s)> = <yj (s), αn xn (s)>,
n=1
<yj (s), f (s)> = αj <yj (s), xj (s)>,
<yj (s), f (s)>
αj = ,
<yj (s), xj (s)>
<yn (s), f (s)>
αn = , n = 1, 2, 3, . . .
<yn (s), xn (s)>
Now in the case at hand, it is easily shown that
1
<yn (s), xn (s)> = , n = 1, 2, 3, . . . ,
2
so we have
αn = 2<yn (s), f (s)>.
The N -term approximate representation of f (s) is thus given by
XN Z 1
f (s) ∼ 2 t/2
e sin (nπt) f (t) dt e−s/2 sin (nπs),
n=1 | 0 | {z }
{z } =xn (s)
=αn
Z 1 N
X
∼ 2 e(t−s)/2 f (t) sin(nπt) sin(nπs) dt,
0 n=1
Z 1 N
X
∼ e(t−s)/2 f (t) (cos(nπ(s − t)) − cos(nπ(s + t))) dt.
0 n=1
1
0.8
0.6
0.4
0.2
s
0.2 0.4 0.6 0.8 1
Figure 7.11: Twenty term Fourier series approximation to a top hat function in terms of a
non-orthogonal basis.
A plot of a twenty-term series expansion of the top hat function is shown in Figure 7.11
In this exercise, the eigenfunctions of the adjoint are actually the reciprocal basis functions. We
see that getting the Fourier coefficients for eigenfunctions of a non-self-adjoint operator requires con-
sideration of the adjoint operator. We also note that it is often a difficult exercise in problems with
practical significance to actually find the adjoint operator and its eigenfunctions.
7.5 Equations
The existence and uniqueness of the solution x of the equation
Lx = y
x = L−1 y (7.73)
2. Eigenvector expansion: Assume that x, y belong to a vector space S and the eigenvec-
tors (e1 , e2 , · · ·) of L span S. Then we can write
X
y = αn en (7.74)
n
X
x = βn en (7.75)
n
where the α’s are known and the β’s are unknown. We get
Lx = y (7.76)
!
X X
L βn en = αn en (7.77)
n n
X X
Lβn en = αn en (7.78)
n n
X X
βn Len = αn en (7.79)
n n
X X
βn λn en = αn en (7.80)
n n
X
(βn λn − αn )en = 0 (7.81)
n
where the λs are the eigenvalues of L. Since the en are linearly independent, we must
demand for all n that
βn λn = αn (7.82)
If all λn 6= 0, then βn = αn /λn and we have the unique solution
X αn
x= en
n
λ n
If, however, one of the λs, λk say, is zero, we still have βn = αn /λn for n 6= k. For
n = k, there are two possibilities:
(a) If αk 6= 0, no solution is possible since equation (7.82) is not satisfied for n = k.
(b) If αk = 0, we have the non-unique solution
X αn
x= en + γek
λn
n6=k
Example 7.43
Solve for x in Lx = y if L = d2 /dt2 , with side conditions x(0) = x(1) = 0, and y(t) = 2t, via an
eigenfunction expansion.
This problem of course has an exact solution via straightforward integration:
d2 x
= 2t; x(0) = x(1) = 0,
dt2
integrates to yield
t 2
x(t) = (t − 1).
3
However, let’s use the series expansion technique. This can be more useful in other problems in
which exact solutions do not exist. First, find the eigenvalues and eigenfunctions of the operator:
d2 x
= λx; x(0) = x(1) = 0.
dt2
This has general solution
√ √
x(t) = A sin −λt + B cos −λt .
This suggests that we expand y(t) = 2t in a Fourier sine series. We know from an earlier problem that
the Fourier sine series for y(t) = 2t is
∞
X 4(−1)n+1
2t = sin(nπt).
n=1
(nπ)
4 1
x(t) ∼ − sin(πt) + 3 sin(2πt),
π3 2π
gives a very good approximation for the solution, which as shown in Figure 7.12, has a peak error of
about 0.008.
x
t error
0.2 0.4 0.6 0.8 1
0.004
-0.02
-0.04 0.002
t
-0.06 0.2 0.4 0.6 0.8 1
-0.08 -0.002
-0.1 -0.004
-0.12 -0.006
-0.008
Figure 7.12: Approximate and exact solution x(t); Error in solution xa (t) − x(t)
Example 7.44
Solve Ax = y using the eigenvector expansion technique when
2 1 3
A= , y= .
1 2 4
Example 7.45
Solve Ax = y using the eigenvector expansion technique when
2 1 3 3
A= , y= , y= .
4 2 4 6
y = c1 e 1 + c2 e 2 .
For this non-symmetric matrix, the eigenvectors are linearly independent, so they form a basis. However
they are not orthogonal, so there is not a direct way to compute c1 and c2 . Matrix inversion shows
that c1 = 52 and c2 = − 21 , so
5 1
y = e1 − e2 .
2 2
Since the eigenvectors form a basis, y can be represented with an eigenvector expansion. However no
λ2 = 0 and c2 6= 0, hence the coefficient β2 = c2 /λ2 does not exist.
solution for x exists because
3
However, for y = , we can say that
6
y = 3e1 + 0e2 .
We note that (3, 6)T is a scalar multiple of the so-called column space vector of A, (2, 4)T . Consequently,
c1 c2
x = e1 + e2 ,
λ1 λ2
c1 0
= e1 + e2 ,
λ1 0
3
= e1 + γe2 ,
4
3 1 −1
= +γ ,
4 2 2
3/4 − γ
= ,
3/2 + 2γ
where γ is an arbitrary constant. Note that the vector e2 = (−1, 2)T lies in the null space of A since
2 1 −1
Ae2 =
4 2 2
0
=
0
Since e2 lies in the null space, any scalar multiple of e2 , say γe2 , also lies in the null space. We can
conclude that for arbitrary y, the inverse does not exist. For vectors y which lie in the column space of
A, the inverse exists, but it is not unique; arbitrary vectors from the null space of A are admitted as
part of the solution.
with certain boundary conditions. Here, L is a differential operator that is not necessarily
linear. We will work with functions and inner products in L2 [a, b] space.
Approximate y(t) by
n
X
y(t) ≈ ya (t) = cj φj (t) (7.88)
j=1
where φj (t), (j = 1, · · · , n) are linearly independent functions (called trial functions) which
satisfy the boundary conditions. Forcing the trial functions to satisfy the boundary condi-
tions, in addition to having æsthetic appeal, makes it much more likely that if convergence
is obtained, the convergence will be to a solution which satisfies the differential equation
and boundary conditions. The trial functions can be orthogonal or non-orthogonal.17 The
constants cj , (j = 1, · · · , n) are to be determined. Substituting into the equation, we get a
residual error
e(t) = Lya (t) − f (t) (7.89)
This error will almost always be non-zero for t ∈ [a, b]. We can, however, choose cj such
that an error, averaged over the domain, is zero. To achieve this, we select now a set of
17
It is occasionally advantageous, especially in the context of what is known as wavelet-based methods, to
add extra functions which are linearly dependent into the set of trial functions. Such a basis is known as a
frame. We will not consider these here; some background is give by Daubechies.
3. Subdomain ψi (t) = 1 for ti−1 ≤ t < ti and zero everywhere else. Note that these
functions are orthogonal to each other. Also this method is easliy shown to reduce to
the well known finite volume method.
∂e
So this method corresponds to ψj = ∂cj
.
5. Moments: ψi (t) = ti , i = 0, 1, · · ·.
If the trial functions are orthogonal and the method is Galerkin, we will, following
Fletcher, who builds on the work of Finlayson, define the method to be a spectral method.
Other less restrictive definitions are in common usage in the present literature, and there is
no single consensus on what precisely constitutes a spectral method.19
18
Boris Gigorievich Galerkin, 1871-1945, Belarussian-born Russian-based engineer and mathematician, a
participant, witness, and victim of much political turbulence, did much of his early great work in the Czar’s
prisons, developed a finite element method in 1915, professor of structural mechanics at St. Petersburg (then
Leningrad).
19
An important school in spectral methods, exemplified in the work of Gottlieb and Orszag, Canuto,
et al., and Fornberg, uses a looser nomenclature, which is not always precisely defined. In these works,
spectral methods are distinguished from finite difference methods and finite element methods in that spectral
methods employ basis functions which have global rather than local support; that is spectral methods’ basis
functions have non-zero values throughout the entire domain. While orthogonality of the basis functions
within a Galerkin framework is often employed, it is not demanded that this be the distinguishing feature
by those authors. Within this school, less emphasis is placed on the framework of the method of weighted
residuals, and the spectral method is divided into subclasses known as Galerkin, tau, and collocation. The
collocation method this school defines is identical to that defined here, and is also called by this school the
Example 7.46
For x ∈ L2 [0, 1], find a one-term approximate solution of the equation
d2 x
+x=t−1
dt2
with x(0) = −1, x(1) = 1.
It is easy to show that the exact solution is
Here we will see how well the method of weighted residuals can approximate this known solution. The
real value of the method is for problems in which exact solutions are not known.
Let y = x − (2t − 1), so that y(0) = y(1) = 0. The transformed differential equation is
d2 y
+ y = −t
dt2
Let us consider a one-term approximation y ≈ ya (t) = cφ(t). There are many choices of basis functions
φ(t). Let’s try finite dimensional non-trivial polynomials which match the boundary conditions. If we
choose φ(t) = a, a constant, we must take a = 0 to satisfy the boundary conditions, so this does not
work. If we choose φ(t) = a + bt, we must take a = 0, b = 0 to satisfy both boundary conditions, so
this also does not work. We can find a quadratic polynomial which is non-trivial and satisfies both
boundary conditions: φ(t) = t(1 − t). Then
We have to determine c. Substituting into the equation, the residual error is found to be
d2 ya
e(t) = Lya − f (t) = + ya − f (t),
dt2
e(t) = −2c + ct(1 − t) − (−t) = t − c(t2 − t + 2).
Then, we choose c such that
Z 1
<ψ(t), e(t)> = <ψ(t), t − c(t2 − t + 2)> = ψ(t) t − c(t2 − t + 2) dt = 0.
0 | {z }
=e(t)
The form of the weighting function ψ is dictated by the particular method we choose:
“pseudospectral” method. In nearly all understandings of the word “spectral,” a convergence rate which is
more rapid than those exhibited by finite difference or finite element methods exists. In fact the accuracy of
a spectral method should grow exponentially with the number of nodes for a spectral method, as opposed
to that for a finite difference or finite element, whose accuracy grows only with the number of nodes raised
to some power.
Another concern which arises with methods of this type is how many terms are necessary to properly
model the desired frequency level. For example, take our equation to be d2 u/dt2 = 1 + u2 ; u(0) = u(π) = 0,
P
and take u = N 2 2
n=1 an sin(nt). If N = 1, we get e(t) = −a1 sin t − 1 − a1 sin t. Expanding the square of the
sin term, we see the error has higher order frequency content: e(t) = −a1 sin t − 1 − a21 (1/2 − 1/2 cos(2t)).
The result is that if we want to get things right at a given level, we may have to reach outside that level.
How far outside we have to reach will be problem dependent.
0.005
0.5
0.0025
t t
0.2 0.4 0.6 0.8 1 0.2 0.4 0.6 0.8 1
-0.0025
-0.5 -0.005
-0.0075
-1 x’’+ x = t - 1; x(0) = -1, x(1) = 1
Figure 7.13: One term estimate xa (t) and exact solution x(t); Error in solution xa (t) − x(t).
1 3c
1. Galerkin: ψ(t) = φ(t) = t(1 − t). The inner product gives 12 − 10 = 0, so that for non-trivial solution,
5
c = 18 = 0.277.
ya (t) = 0.277t(1 − t).
xa (t) = 0.277t(1 − t) + 2t − 1.
xa (t) = 0.286t(1 − t) + 2t − 1.
xa (t) = 0.273t(1 − t) + 2t − 1.
∂e(t) 11 101 55
4. Least squares: ψ(t) = ∂c = −t2 + t − 2. Thus − 12 + 30 c = 0, from which c = 202 = 0.273.
xa (t) = 0.273t(1 − t) + 2t − 1.
5. Moments: ψ(t) = 1 which, for this case, is the same as the subdomain method above.
xa (t) = 0.273t(1 − t) + 2t − 1.
The approximate solution determined by the Galerkin method is overlaid against the exact solution in
Figure 7.13. Also shown is the error in the approximation. The approximation is surprisingly accurate.
Some simplification can arise through use of integration by parts. This has the result of
admitting basis functions which have less stringent requirements on the continuity of their
derivatives. It is also a commonly used strategy in the finite element technique.
Example 7.47
Consider a slight variant of the previous example problem, and employ integration by parts.
d2 y
+ y = f (t), y(0) = 0, y(1) = 0.
dt2
Again, take a one term expansion
ya (t) = cφ(t).
At this point, we will only require φ(t) to satisfy the boundary conditions, and will specify it later. The
error in the approximation is
d2 ya d2 φ
e(t) = + y a − f (t) = c + cφ − f (t).
dt2 dt2
Now set a weighted error to zero. We will also require the weighting function ψ(t) to vanish at the
boundaries. Z 1 2
d φ
<ψ, e> = ψ(t) c 2 + cφ(t) − f (t) dt = 0.
0 dt
| {z }
=e(t)
Rearranging, we get
Z 1 Z 1
d2 φ
c ψ(t) 2 + ψ(t)φ(t) dt = ψ(t)f (t) dt.
0 dt 0
So, the basis function φ only needs an integrable first derivative rather than an integrable second
derivative. As an aside, we note that the term on the left hand side bears resemblance (but differs by
1
a sign) to an inner product in the Sobolov space W2 [0, 1] in which the Sobolov inner product <., .>s
R1 dψ dφ
(an extension of the inner product for Hilbert space) is <ψ(t), φ(t)>s = 0 ψ(t)φ(t) + dt dt dt.
Taking now, as before, φ = t(1 − t) and then choosing a Galerkin method so ψ(t) = φ(t) = t(1 − t),
and f (t) = −t, we get
Z 1 Z 1
2 2 2
c t (1 − t) − (1 − 2t) dt = t(1 − t)(−t) dt,
0 0
which gives
3 1
c − =− ,
10 12
so
5
c= ,
18
as was found earlier. So
5
ya = t(1 − t),
18
with the Galerkin method.
Example 7.48
For y ∈ L2 [0, 1], find a two-term spectral approximation (which by our definition of “spectral”
mandates a Galerkin formulation) to the solution of
d2 y √
+ t y = 1; y(0) = 0, y(1) = 0.
dt2
Let’s try polynomial basis functions. At a minimum, these basis functions must satisfy the boundary
conditions. Assumption of the first basis function to be a constant or linear gives rise to a trivial basis
function when the boundary conditions are enforced. The first non-trivial basis function is a quadratic:
φ1 (t) = a0 + a1 t + a2 t2
We need φ1 (0) = 0 and φ1 (1) = 0. The first condition gives a0 = 0; the second gives a1 = −a2 , so we
have φ1 = a1 (t − t2 ). Since the magnitude of a basis function is arbitrary, a1 can be set to unity to give
Alternatively, we could have chosen the magnitude in such a fashion to guarantee an orthonormal basis
function, but that is a secondary concern for the purposes of this example.
We need a second linearly independent basis function for the two term approximation. We try a
third order polynomial:
φ2 (t) = b0 + b1 t + b2 t2 + b3 t3 .
Enforcing the boundary conditions as before gives b0 = 0 and b1 = −(b2 + b3 ), so
φ2 (t) = −(b2 + b3 )t + b2 t2 + b3 t3 .
To achieve a spectral method (which in general is not necessary to achieve an approximate solution!),
we enforce <φ1 , φ2 > = 0:
Z 1
t(1 − t) −(b2 + b3 )t + b2 t2 + b3 t3 dt = 0
0 | {z } | {z }
=φ1 (t) =φ2 (t)
b2 b3
− − = 0
30 20
3
b2 = − b3
2
Substituting and factoring gives
b3
φ2 (t) = t(1 − t)(2t − 1).
2
Again, because φ2 is a basis function, the lead constant is arbitrary; we take for convenience b3 = 2 to
give
φ2 = t(1 − t)(2t − 1).
Again, b3 could alternatively have been chosen to yield an orthonormal basis function.
Now we want to choose c1 and c2 so that our approximate solution
e(t) = Lya (t) − f (t) = L (c1 φ1 (t) + c2 φ2 (t)) − 1 = c1 Lφ1 (t) + c2 Lφ2 (t) − 1.
<ψ1 , e> = c1 <ψ1 , Lφ1 > + c2 <ψ1 , Lφ2 > − <ψ1 , 1> = 0
<ψ2 , e> = c1 <ψ2 , Lφ1 > + c2 <ψ2 , Lφ2 > − <ψ2 , 1> = 0
This is easily cast in matrix form as a linear system of equations for the unknowns c1 and c2
<ψ1 , Lφ1 > <ψ1 , Lφ2 > c1 <ψ1 , 1>
=
<ψ2 , Lφ1 > <ψ2 , Lφ2 > c2 <ψ2 , 1>
Each of the inner products represents a definite integral which is easily evaluated via computer algebra.
For example,
Z 1 215
<φ1 , Lφ1 > = t(1 − t) −2 + (1 − t)t3/2 dt = − .
0 693
When each inner product is evaluated, the following system results
215 16
1
− 693 9009 c1 6
= .
16 197
9009 − 1001 c2 0
The two-term approximate solution determined is overlaid against a more accurate solution obtained by
numerical integration of the full equation in Figure 7.14. Also shown is the error in the approximation.
The two term solution is surprisingly accurate.
By normalizing the basis functions, we can find an orthonormal expansion. One finds that
s s
Z 1 Z 1
1
||φ1 ||2 = φ21 dt = t2 (1 − t)2 dt = √
0 0 30
s s
Z 1 Z 1
2 1
||φ2 ||2 = φ2 dt = t2 (1 − t)2 (2t − 1)2 dt = √
0 0 210
-0.06
0.0005
-0.08
Figure 7.14: Two term spectral (Galerkin) estimate ya (t) and highly accurate numerical
solution y(t); Error in approximation ya (t) − y(t).
Example 7.49
For the equation of the previous example,
d2 y √
+ t y = 1; y(0) = 0; y(1) = 1,
dt2
examine the convergence rates for a collocation method as the number of modes becomes large.
Let us consider a set of trial functions which do not happen to be orthogonal, but are, of course,
linearly independent. Take
φi (t) = ti (t − 1), i = 1, . . . , n.
So we seek to find a vector c = ci , i = 1, . . . , n, such that for a given number of collocation points n the
approximation
yn (t) = c1 φ1 (t) + . . . ci φi (t) + . . . + cn φn (t),
drives a weighted error to zero. Obviously each these trial functions satisfies both boundary conditions,
and they have the advantage of being easy to program for an arbitrary number of modes, as no Gram-
Schmidt orthogonalization process is necessary. The details of the analysis are similar to those of the
previous example, except we perform it many times, varying the number of nodes in each calculation.
For the collocation method, we take the weighting functions to be
ψj (t) = δ(t − tj ), j = 1, . . . , n.
||yn - y ||2
N
10 - 3
10 - 4
10 - 5 - 21.9
||yn - y ||2 ~ n
N
-6
10
1 2 3 5 7 10 n
Figure 7.15: Error in solution yn (t) − yN (t) as a function of number of collocation points n.
Here we choose tj = j/(n + 1), j = 1, . . . , n, so that the collocation points are evenly distributed in
t ∈ [0, 1]. We then form the matrix
<ψ1 , Lφ1 >, <ψ1 , Lφ2 > ... <ψ1 , Lφn >
<ψ2 , Lφ1 >, <ψ2 , Lφ2 > ... <ψ2 , Lφn >
A=
.. .. .. .. ,
. . . .
<ψn , Lφ1 >, <ψn , Lφ2 > . . . <ψn , Lφn >
A plot of the error en is plotted as a function of n in Figure 7.15. We notice even on a logarithmic plot
that the error reduction is accelerating as the number of nodes n increases. If the slope had relaxed
to a constant, then the convergence would be a power law convergence; which is characteristic of finite
difference and finite element methods. For this example of the method of weighted residuals, we see that
the rate of convergence increases as the number of nodes increases, which is characteristic of exponential
convergence. For exponential convergence, we have en ∼ exp(−αn), where α is some positive constant;
for power law convergence, we have en ∼ n−β where β is some positive constant. At the highest value
of n, n = 10, we have a local convergence rate of O(n−21.9 ) which is remarkably fast. In comparison,
a second order finite difference technique will converge at a rate of O(n−2 ). In general and if possible
one would choose a method with the fastest convergence rate, all else being equal.
Problems
1. Use a one-term collocation method with a polynomial basis function to find an approximation for
y ′′′′ + (1 + x)y = 1
n
!2 n
! n
!
X X X
xi yi ≤ x2i yi2
i=1 i=1 i=1
4. If x, y ∈ X, an inner product space, and x is orthogonal to y, then show that ||x + αy|| = ||x − αy||
where α is a scalar.
5. For an inner product space, show that
(a) A is the set of all polynomials which are all exactly of degree n, F = R.
(b) A is the set of all functions with continuous second derivatives over the interval [0, L] and
satisfying the differential equation y ′′ + 2y ′ + y = 0, F = R.
(c) A = R, F = R.
9. Given a set S of linearly independent vectors in a vector space V, show that any subset of S is also
linearly independent.
10. Do the following vectors, (3, 1, 4, −1)T , (1, −4, 0, 4)T , (−1, 2, 2, 1)T , (−1, 9, 5, −6)T , form a basis in R4 ?
11. Given x1 , the iterative procedure xn+1 = Txn generates x2 , x3 , x4 , · · ·, where T is a linear operator
and all the x’s belong to a complete normed space. Show that {xn , n = 1, 2, · · ·} is a Cauchy sequence
if ||T|| < 1. Does it converge? If so find the limit.
12. If {ePk , k = 1, 2, · · ·} is an orthonormal set in a Hilbert space H, show that for every x ∈ H, the vector
n
y = k=1 <x, ek >ek exists in H, and that x − y is orthogonal to every ek .
2 −4
13. Let the linear operator A : C2 → C2 be represented by the matrix A = . Find ||A|| if all
1 5
vectors in the domain and range are within a Hilbert space.
2 2 2 + i −4
14. Let the linear operator A : C → C be represented by the matrix A = . Find ||A||
1 5
if all vectors in the domain and range are within a Hilbert space.
Rb
15. Using the inner product (x, y) = a w(t)x(t)y(t) dt, where w(t) > 0 for a ≤ t ≤ b, show that the
Sturm-Liouville operator
1 d d
L= p(t) + r(t)
w(t) dt dt
with αx(a) + βx′ (a) = 0, and γx(b) + δx′ (b) = 0 is self-adjoint.
16. For elements x, y and z of an inner product space, prove the Appolonius identity:
2
1 1
||z − x||22 + ||z − y||22 2
= ||x − y||2 + 2 z − (x + y)
2 2 2
17. If x, y ∈ X an inner product space, and x is orthogonal to y, then show that ||x + ay||2 = ||x − ay||2
where a is a scalar.
18. Using the Gram-Schmidt procedure, find the first three members of the orthonormal set belonging to
L2 (−∞, ∞), using the basis functions {exp(−t2 /2), t exp(−t2 /2), t2 exp(−t2 /2), · · ·}. You may need
the following definite integral
Z ∞ √
exp(−t2 /2) dt = 2π.
−∞
19. Let C(0,1) be the space of all continuous functions in (0,1) with the norm
s
Z 1
||f ||2 = |f (t)|2 dt.
0
Show that
2n tn+1 for 0 ≤ t < 12
fn (t) =
1 − 2n (1 − t)n+1 for 12 ≤ t ≤ 1
belongs to C(0,1). Show also that {fn , n = 1, · · ·} is a Cauchy sequence, and that C(0,1) is not
complete.
20. Find the first three terms of the Fourier-Legendre series for f (x) = cos(πx/2) for x ∈ [−1, 1]. Compare
graphically with exact function.
21. Find the first three terms of the Fourier-Legendre series for
−1 for − 1 ≤ x < 0
f (x) =
1 for 0 ≤ x ≤ 1
22. Consider
d3 y
+ 2t3 y = 1 − t
dt3
dy
y(0) = 0 y(2) = 0 (0) = 0
dt
Choosing polynomials as the basis functions, use a Galerkin and moments method to obtain a two-
term estimate to y(t). Plot your approximations and the exact solution on a single curve. Plot the
error in both methods for x ∈ [0, 2]
23. Solve
x′′ + 2xx′ + t = 0
with x(0) = 0, x(4) = 0, approximately using a two-term weighted residual method where the basis
functions are of the type sin λt. Do both a spectral (as a consequence Galerkin) and pseudospectral
(as a consequence collocation) method. Plot your approximations and the exact solution on a single
curve. Plot the error in both methods for x ∈ [0, 4].
x1 + 3x2 + x3 − x4 = 0
−2x1 + 2x2 − x3 + x4 = 0
form a vector space. Find the dimension and a set of basis vectors.
25. Let
1 1 1
A= 0 1 1
0 0 1
26. For any complete orthonormal set {φi , i = 1, 2, · · ·} in a Hilbert space H, show that
X
u = <u, φi >φi
i
X
<u, v> = <u, φi ><v, φi >
i
X
||u||22 = |<u, φi >|2
i
d2 y
+ y4 = 1
dx2
with y(0) = 0, y(1) = 0.
31. Show that s s s
Z b Z b Z b
2 2 2
(f (x) + g(x)) dx ≤ (f (x)) dx + (g(x)) dx
a a a
d2 x
+ k 2 x = f (t), with x(0) = a, x(π) = b
dt2
Write the solution of the differential equation in terms of this function.
38. Find the first three terms of the Fourier-Legendre series for
−2 for − 1 ≤ x < 0
f (x) =
1 for 0 ≤ x ≤ 1
43. Consider functions of two variables in a domain Ω with the inner product defined as
ZZ
<u, v> = u(x, y)v(x, y) dx dy
Ω
Find the space of functions such that the Laplacian operator is self-adjoint.
44. Find the eigenvalues and eigenfunctions of the operator L where
d2 y dy
Ly = (1 − t2 ) −t
dt2 dt
with t ∈ [−1, 1] and y(−1) = y(1) = 0. Show that there exists a weight function r(x) such that the
eigenfunctions are orthogonal in [−1, 1] with respect to it.
45. Show that the eigenvalues of an operator and its adjoint are complex conjugates of each other.
46. Using an eigenvector expansion, find the general solution of A · x = y where
2 0 0
A = 0 1 1
0 1 1
2
y = 3
5
47. Show graphically that the Fourier trigonometric series representation of the function
−1 if −π ≤ t < 0
f (t) =
1 if 0 ≤ t ≤ π
always has an overshoot near x = 0, however many terms one takes (Gibbs phenomenon). Estimate
the overshoot.
48. Let {e1 , · · · , en } be an orthonormal set in an inner product space S. Approximate x ∈ S by y =
β1 e1 + · · · + βn en , where the β’s are to be selected. Show that ||x − y|| is a minimum if we choose
βi = <x, ei >.
49. (a) Starting with a vector in the direction (1, 2, 0)T use the Gram-Schmidt procedure to find a set of
orthonormal vectors in R3 . Using these vectors, construct (b) an orthogonal matrix Q, and then find
(c) the angles between xi and Qxi , where xi is (1, 0, 0)T , (0, 1, 0)T and (0, 0, 1)T respectively. The
orthogonal matrix Q is defined as a matrix having orthonormal vectors in its columns.
2
d
50. Find the null space of the operator L defined by Lx = dt 2 x(t). Also find the eigenvalues and eigen-
dx
functions (in terms of real functions) of L with x(0) = 1, dt (0) = 0.
51. Find all approximate solutions of the boundary value problem
d2 y
+ y + 5y 2 = −x,
dx2
with y(0) = y(1) = 0 using a two term collocation method. Compare graphically with the exact
solution determined by numerical methods.
52. Find a one-term approximation for the boundary value problem
y ′′ − y = −x3
with y(0) = y(1) = 0, using the collocation, Galerkin, least-squares, and moments methods. Compare
graphically with the exact solution.
1+ 1
53. Consider the sequence { 2+ n1 } in Rn . Show that this is a Cauchy sequence. Does it converge?
n
54. Prove that (Ta Tb )∗ = T∗b T∗a when Ta and Tb are linear operators which operate on vectors in a
Hilbert space.
55. If {xi } is a sequence in an inner product space such that the series ||x1 || + ||x2 || + · · · converges, show
that {sn } is a Cauchy sequence, where sn = x1 + x2 + · · · + xn .
2
56. If L(x) = a0 (t) ddt2x + a1 (t) dx
dt + a2 (t)x, find the operator that is formally adjoint to it.
57. If Z t
y(t) = A[x(t)] = x(τ ) dτ
0
where y(t) and x(t) are real functions in some properly defined space, find the eigenvalues and eigen-
functions of the operator A.
58. Using a dual basis, expand the vector (1, 3, 2)T in terms of the basis vectors (1, 1, 1)T , (1, 0, −1)T , and
(1, 0, 1)T in R3 . The inner product is defined as usual.
59. With f1 (x) = 1 + i + x and f2 (x) = 1 + ix + ix2
a) Find the L2 [0, 1] norms of f1 (x) and f2 (x).
b) Find the inner product of f1 (x) and f2 (x) under the L2 [0, 1] norm.
c) Find the “distance” between f1 (x) and f2 (x) under the L2 [0, 1] norm.
60. Show the vectors u1 = (−i, 0, 2, 1 + i)T , u2 = (1, 2, i, 3)T , u3 = (3 + i, 3 − i, 0, −2)T , u4 = (1, 0, 1, 3)T
form a basis in C4 . Find the set of reciprocal basis vectors. For x ∈ C4 , and x = (i, 3 − i, −2, 2)T ,
express x as an expansion in the above defined basis vectors. That is find ci such that x = ci ui
61. The following norms can be used in Rn , where x = (ξ1 , · · · , ξn ) ∈ Rn .
(a) ||x||∞ = max1≤j≤n |ξj |
Pn
(b) ||x||1 = j=1 |ξj |
Pn
(c) ||x||2 = ( j=1 |ξj |2 )1/2
Pn
(d) ||x||p = ( j=1 |ξj |p )1/p , 1 ≤ p < ∞
Show by examples that these are all valid norms.
62. Show that the set of all matrices A : Rn → Rn is a vector space under the usual rules of matrix
manipulation.
63. Show that if A is a linear operator such that
Pn
(a) A : (Rn , || · ||∞ ) → (Rn , || · ||1 ) then ||A|| = i,j=1 Aij .
n n
Pn
(b) A : (R , || · ||∞ ) → (R , || · ||∞ ) then ||A|| = max1≤i≤n j=1 Aij .
Linear algebra
A · x = b, (8.1)
where n, m ∈ N are the positive integers which give the dimensions. If n = m, the matrix
is square, and solution techniques are usually straightforward. For n 6= m, which arises
often in physical problems, the issues are not as straightforward. In some cases we find an
infinite number of solutions; in others we find none. Relaxing our equality constraint, we
can, however, always find a vector x∗
This vector x∗ is the best solution to the equation A · x = b, for cases in which there is no
exact solution. Depending on the problem, it may turn out that x∗ is not unique. It will
always be the case, however, that of all the vectors x∗ which minimize ||A · x − b||2 , that
one of them, x̂, will itself have a minimum norm.
287
288 CHAPTER 8. LINEAR ALGEBRA
• det An×n is equal to the volume of a parallelepiped in n-dimensional space whose edges
are formed by the rows of A.
• If all elements of a row (or column) are multiplied by a scalar, the determinant is also
similarly multiplied.
• The elementary operation of subtracting a multiple of one row from another leaves the
determinant unchanged.
• If two rows (or columns) of a matrix are interchanged the sign of the determinant
changes.
A singular matrix is one whose determinant is zero. The rank of a matrix is the size r of
the largest square non-singular matrix that can be formed by deleting rows and columns.
While the determinant is useful to some ends in linear algebra, most of the common
problems are better solved without using the determinant at all; in fact it is probably a fair
generalization to say that the determinant is less, rather than more, useful than imagined by
many. It is useful in solving linear systems of equations of small dimension, but becomes much
too cumbersome relative to other methods for commonly encountered large systems of linear
algebraic equations. While it can be used to find the rank, there are also other more efficient
means to calculate this. Further, while a zero value for the determinant almost always has
significance, other values do not. Some matrices which are particularly ill-conditioned for
certain problems often have a determinant which gives no clue as to difficulties which may
arise.
where the elements of C are obtained by adding the corresponding elements of A and B.
Multiplication of a matrix by a scalar α can be defined as
αAn×m = Bn×m
Also
• Any vector x ∈ Cm can be written as a linear combination of vectors in the row space
and the right null space.
• Any m dimensional vector x which is in the right null space of A is orthogonal to any
m dimensional vector in the row space. This comes directly from the definition of the
right null space A · x = 0.
• Any vector y ∈ Cn can be written as the sum of vectors in the column space and the
left null space.
• Any n dimensional vector y which is in the left null space of A is orthogonal to any
n dimensional vector in the column space. This comes directly from the definition of
the left null space yT · A = 0.
Example 8.1
Find the column and row spaces of
1 0 1
A=
0 1 2
The column space consists of the vectors α1 c1 + α2 c2 + α3 c3 , where the α’s are any scalars. Since only
two of the ci ’s are linearly independent, the dimension of the column space is also two. We can see this
by looking at the subdeterminant
1 0
det = 1,
0 1
which indicates the rank, r = 2.
Note that
• c1 + 2c2 = c3
• The three column vectors thus lie in a single two-dimensional plane.
• The three column vectors are thus said to span a two-dimensional subspace of R3
The two row vectors are
r1 = 1 0 1
r2 = 0 1 2
The row space consists of the vectors β1 r1 + β2 r2 , where the β’s are any scalars. Since the two ri ’s are
linearly independent, the dimension of the row space is also two. That is the two row vectors are both
three dimensional, but span a two-dimensional subspace.
We note for instance, if x = (1, 2, 1)T , that A · x = b gives
1
1 0 1 2
2 = .
0 1 2 4
1
So
b = 1c1 + 2c2 + 1c3 .
That is b is a linear combination of the column space vectors and thus lies in the column space of A.
We note for this problem that since an arbitrary b is two-dimensional and the dimension of the column
space is two, that we can represent an arbitrary b as some linear combination of the column space
vectors. For example, we can also say that b = 2c1 + 4c2 . We also note that x in general does not
lie in the row space of A, since x is an arbitrary three-dimensional vector, and we only have enough
row vectors to span a two-dimensional subspace (i.e. a plane embedded in a three-dimensional space).
However, as will be seen, x does lie in the space defined by the combination of the row space of A, and
the right null space of A (the set of vectors x for which A · x = 0). In special cases, x will in fact lie
in the row space of A.
It may be better to say here that A is a linear operator which operates on elements which are
in a space of dimension k × m so as to generate elements which are in a space of dimension
n × m; that is, A : Rk × Rm → Rn × Rm .
Example 8.2
Consider the matrix operator
1 2 1
A=
−3 3 1
which operates on 3 × 4 matrices, i.e.
A : R3 × R4 → R2 × R4
A vector operating on a vector can yield a scalar or a matrix, depending on the order of
operation.
Example 8.3
Consider the vector operations A1×3 · B3×1 and B3×1 · A1×3 where
A1×3 = aT = ( 2 3 1)
3
B3×1 = b = −2
5
Then
3
A1×3 · B3×1 = aT · b = ( 2 3 1 ) −2 = (2)(3) + (3)(−2) + (1)(5) = 5
5
This is the ordinary inner product <a, b>. The commutation of this operation however yields a matrix:
3 (3)(2) (3)(3) (3)(1) 6 9 3
B3×1 · A1×3 = baT = −2 ( 2 3 1 ) = (−2)(2) (−2)(3) (−2)(1) = −4 −6 −2
5 (5)(2) (5)(3) (5)(1) 10 15 5
This is the dyadic product of the two vectors. Note that for vector (lower case notation) the dyadic
product usually is not characterized by the “dot” operator that we use for the vector inner product.
A special case is that of a square matrix An×n of size n. For square matrices of the same
size both A · B and B · A exist. While A · B and B · A both yield n × n matrices, the
actual value of the two products is different. In what follows, we will often assume that we
are dealing with square matrices.
Properties of matrices include
1. (A · B) · C = A · (B · C) (associative),
2. A · (B + C) = A · B + A · C (distributive),
3. (A + B) · C = A · C + B · C (distributive),
where the unsubscripted identity matrix is understood to be square with the correct dimen-
sion for matrix multiplication.
The transpose AT of a matrix A is one in which the terms above and below the diagonal
are interchanged. For any matrix An×m , we find that A · AT and AT · A are square matrices
of size n and m, respectively.
Properties include
1. det A = det AT ,
2. (An×m · Bm×n )T = BT · AT ,
1 1
A = (A + AT ) + (A − AT ), (8.10)
2 2
where 21 (A + AT ) is symmetric and 21 (A − AT ) is anti-symmetric.
A lower (or upper) triangular matrix is one in which all entries above (or below) the main
diagonal are zero. Lower triangular matrices are often denoted by L, and upper triangular
matrices by either U or R.
A positive definite matrix A is a matrix for which xT · A · x > 0 for all nonzero vectors x.
A positive definite matrix has real, positive eigenvalues. There exists a nonsingular W such
that A = WT · W. All the eigenvalues of such a matrix are positive. Every positive definite
matrix A can be written as A = L · LT , where L is a lower triangular matrix (Cholesky
decomposition).
A permutation matrix P is a square matrix composed of zeroes and a single one in each
column. None of the ones occur in the same row. It effects a row exchange when it operates
on a general matrix A. It is never singular, and is in fact its own inverse, P = P−1 , so
P · P = I. Also ||P||2 = 1.
Example 8.4
Find P which effects the exchange of the first and second rows of A, where
1 3 5 7
A = 2 3 1 2.
3 1 3 2
To construct P, we begin with at 3 × 3 identity matrix I. For a first and second row exchange, we
replace the ones in the (1, 1) and (2, 2) slot with zero, then replace the zeroes in the (1, 2) and (2, 1)
slot with ones. Thus
0 1 0 1 3 5 7 2 3 1 2
P · A = 1 0 02 3 1 2 = 1 3 5 7.
0 0 1 3 1 3 2 3 1 3 2
Example 8.5
Find the rank and right null space of
1 0 1
A= 5 4 9 .
2 4 6
x1 + x3 = 0,
5x1 + 4x2 + 9x3 = 0,
2x1 + 4x2 + 6x3 = 0.
One strategy to solve singular systems is to take one of the variables to be a known parameter, and see
if the resulting system can be solved. If the resulting system remains singular, take a second variable
to be a second parameter. This ad hoc method will later be made systematic.
So here take x1 = t, and consider the first two equations, which gives
0 1 x2 −t
= .
4 9 x3 −5t
Solving, we find x2 = t, x3 = −t. So,
x1 t 1
x = x2 = t = t 1 , t ∈ R1 .
x3 −t −1
Therefore, the right null space is the straight line in R3 which passes through (0,0,0) and (1,1,-1).
8.2.3.2 Inverse
(−1)i+j bji
a−1
ij = , (8.11)
det A
where bij is the minor of aji which is the determinant of the matrix obtained by canceling
out the j-th row and i-th column. The inverse of a diagonal matrix is also diagonal, but
with the reciprocals of the original diagonal elements.
Example 8.6
Find the inverse of
1 1
A= .
−1 1
The inverse is
−1 1 1 −1
A = .
2 1 1
We can confirm that A · A−1 = A−1 · A = I.
8.2.4 Equations
In general, for matrices that are not necessarily square, the equation An×m · xm×1 = bn×1
is solvable iff b can be expressed as combinations of the columns of A. Problems in which
m < n are overconstrained; in special cases, those in which b is in the column space of A,
a unique solution x exists. However in general no solution x exists; nevertheless, one can
find an x which will minimize ||A · x − b||2 . This is closely related to what is known as
the method of least squares. Problems in which m > n are generally underconstrained, and
have an infinite number of solutions x which will satisfy the original equation. Problems for
which m = n (square matrices) have a unique solution x when the rank r of A is equal to
n. If r < n, then the problem is underconstrained.
Example 8.7
For x ∈ R2 , b ∈ R3 , consider A : R2 → R3 ,
1 2 5
1 x1
0 = 1.
x2
1 1 3
The column space of A is spanned by the two column vectors
1 2
c1 = 1 , c2 = 0 .
1 1
Our equation can also be cast in the form which makes the contribution of the column vectors obvious:
1 2 5
x1 1 + x2 0 = 1 .
1 1 3
Here we have the unusual case that b = (5, 1, 3)T is in the column space of A (in fact b = c1 + 2c2 ),
and we have a unique solution of
1
x= .
2
Note that the solution vector x lies in the row space of A; here it identically the first row vector
r1 = (1, 2)T . Note also that here the column space is a two-dimensional subspace, in this case a plane
defined by the two column vectors, embedded within a three-dimensional space. The operator A maps
arbitrary two-dimensional vectors x into the three-dimensional b; however, these b vectors are confined
to a two-dimensional subspace within the greater three-dimensional space. Consequently, we cannot
always expect to find a vector x for arbitrary b!
3
R
c2
2 A.x = b
C c1
Figure 8.1: Plot for b which lies in column space (space spanned by c1 and c2 ) of A.
A sketch of this system is shown in Figure 8.1. Here we sketch what might represent this example
in which the column space of A does not span the entire space R3 , but for which b lies in the column
space of A. In such a case ||A · x − b||2 = 0. We have A as a matrix which maps two dimensional
vectors x into three dimensional vectors b. Our space is R3 and embedded within that space are two
column vectors c1 and c2 which span a column space C2 , which is represented by a plane within a three
dimensional volume. Since b in this example happens to lie in the column space, there exists a unique
vector x for which A · x = b.
Example 8.8
Consider now
1 2 0
1 0 x1
= 1.
x2
1 1 3
Here it turns out that b = (0, 1, 3)T is not in the column space of A, and there is no solution x for
which A · x = b! Again, the column space is a plane defined by two vectors; the vector b does not
happen to lie in the plane defined by the column space. However, we can find a solution x = xp , where
xp can be shown to minimize the least squares error e = ||A·xp − b||2 . This is achieved by the following
procedure in which we operate on both vectors A · xp and b by the operator AT so as to map both
vectors into the same space, namely the row space of A. Once the vectors are in the same space, a
unique inversion is possible.
A · xp ≃ b
T
A · A · xp = AT · b
= (AT · A)−1 · AT · b
xp
1 2 0
1 1 1 x1 1 1 1
1 0 = 1
2 0 1 x2 2 0 1
1 1 3
3 3 x1 4
=
3 5 x2 3
11
x1 6
=
x2 − 21
3
R
c2
2 . = bp = b
C c 1A x p
Figure 8.2: Plot for b which lies outside of column space (space spanned by c1 and c2 ) of A.
Note the resulting xp will not satisfy A · xp = b. In fact, ||A · xp − b||2 = 2.0412. If we tried any nearby
T
x, say x = 2, − 53 , ||A · x − b||2 = 2.0494 > 2.0412. Since the problem is linear, this minimum is
global; if we take x = (10, −24)T , then ||A · x − b||2 = 42.5911 > 2.0412. Though we have not proved
it, our xp is the unique vector which minimizes the least squares error.
Further manipulation shows that we can write our solution as a combination of vectors in the row
space of A. As the dimension of the right null space of A is zero, there is no possible contribution from
the right null space vectors. 11
6 1 1
= α1 + α2 .
− 12 2 0
11
6 1 1 α1
= .
− 12 2 0 α2
Solving, we find
α1 − 41
= 25 .
α2 12
So
x1 1 1 25 1
= − + .
x2 4 2 12 0
| {z }
linear combination of row space vectors
We could also have chosen to expand in terms of the other row space vector (1, 1)T , since any two of
the three row space vectors span the space R2 .
The vector A · xp actually represents the projection of b onto the subspace spanned by the column
vectors (i.e. the column space). Call the projected vector bp :
bp = A · xp = A · (AT · A)−1 · AT ·b
| {z }
projection matrix
4 T
For this example bp = 56 , 11
6 , 3 . We can think of bp as the shadow cast by b onto the column space.
A sketch of this system is shown in in Figure 8.2 Here we sketch what might represent this example
in which the column space of A does not span the entire space R3 , and for which b lies outside of
the column space of A. In such a case ||A · xp − b||2 > 0. We have A as a matrix which maps two
dimensional vectors x into three dimensional vectors b. Our space is R3 , and that embedded within
that space are two column vectors c1 and c2 which span a column space C2 , which is represented by a
plane within a three dimensional volume. Since b lies outside the column space, there exists no unique
vector x for which A · x = b.
Example 8.9
Consider now A : R3 → R2 such that
x
1 1 1 1 1
x2 = .
2 0 1 3
x3
Certainly b = (1, 3)T lies in the column space of A, since for example, b = 0(1, 2)T − 2(1, 0)T + 3(1, 1)T .
Setting x1 = t, where t is an arbitrary number, lets us solve for x2 , x3 :
t
1 1 1 1
x2 = ,
2 0 1 3
x3
1 1 x2 1−t
= .
0 1 x3 3 − 2t
Inversion gives
x2−2 + t
= ,
x3 3 − 2t
so
x1 t 0 1
x2 = −2 + t = −2 + t 1 , t ∈ R1 .
x3 3 − 2t 3 −2
| {z }
right null space
A useful way to think of problems such as this which are underdetermined is that the matrix A maps
the additive combination of a unique vector from the row space of A plus an arbitrary vector from the
right null space of A into the vector b. Here the vector (1, 1, −2)T is in the right null space; however,
the vector (0, −2, 3)T has components in both the right null space and the row space. Let us extract
the parts of (0, −2, 3)T which are in each space. Since the row space and right null space are linearly
indpendent, they form a basis, and we can say
0 1 2 1
−2 = a1 1 + a2 0 + a3 1 .
3 1 1 −2
| {z } | {z }
row space right null space
So x can be rewritten as
1 2 1
2 4
x=− 1 + 0 + t− 1 , t ∈ R1 .
3 3
1 1 −2
| {z } | {z }
row space right null space
The first two terms in the final expression above are the unique linear combination of the row space
vectors, while the third term is from the right null space. As by definition, A maps any vector from the
right null space into the zero element, it makes no contribution to forming b; hence, one can allow for
an arbitrary constant. Note the analogy here with solutions to inhomogeneous differential equations.
The right null space vector can be thought of as a solution to the homogeneous equation, and the terms
with the row space vectors can be thought of as particular solutions.
We can also write the solution x in matrix form. The matrix is composed of three column vectors,
which are the original two row space vectors and the right null space vector, which together form a
basis in R3 :
2
1 2 1 −3
x = 1 0 1 1 , t ∈ R1 .
4
1 1 −2 t− 3
While the right null space vector is orthogonal to both row space vectors, the row space vectors are not
orthogonal to themselves, so this basis is not orthogonal. Leaving out the calculational details, we can
use the Gram-Schmidt procedure to cast the solution on an orthonormal basis:
√1
√1 1
−√ 6
√ 12 √
3
1 √1
4 √1
t ∈ R1 .
x= √ 3 − 2 √ + 6 t− q6 ,
3 2 3
√1 0 − 2
3 3
| {z } | {z }
row space right null space
The first two terms are in the row space, now represented on an orthonormal basis, the third is in the
right null space. In matrix form, we can say that
√1 − √12 √1
3 6 √1
√1 √1 √1 √3
t ∈ R1 .
x= 3 2 q6 √ − 2 ,
√1 0 − 2 6 t − 34
3 3
Of course, there are other orthonormal bases on which the system can be cast.
We see that the minimum length of the vector x occurs when t = 43 , that is when x is entirely in
the row space. In such a case we have
s 2 r
1 √ 2 7
min||x||2 = √ + − 2 =+ .
3 3
Lastly note that here, we achieved a reasonable answer by setting x1 = t at the outset. We could
have achieved an equivalent result by starting with x2 = t, or x3 = t. This will not work in all problems,
as will be discussed in the section on row echelon form.
Example 8.10
Consider A : R4 → R3 such that
x1
1 2 0 4 1
3 2 −1 3 x2
= 3
x3
−1 2 1 5 2
x4
Using elementary row operations to perform Gaussian elimination gives rise to the equivalent system:
x1
1 0 −1/2 −1/2 0
x
0 1 1/4 9/4 2 = 0
x3
0 0 0 0 1
x4
We immediately see that there is a problem in the last equation, which purports 0 = 1! What is actually
happening is that A is not full rank r = 3, but actually has r = 2, so vectors x ∈ R4 are mapped into a
two-dimensional subspace. So, we do not expect to find any solution to this problem, since our vector b
is an arbitrary three dimensional vector which most likely does not lie in the two-dimensional subspace.
We can, however, find an x which minimizes the least squares error. We return to the original equation
and operate on a both sides with AT to form AT · A · x = AT · b. It can be easily verified that if we
chose to operate on the system which was reduced by Gaussian elimination that we would not recover
a solution which minimized ||A · x − b||!
x1 1 3 −1
1 3 −1
1 2 0 4 1
2 2 2 x 2 2 2
3 2 −1 3 2 = 3
0 −1 1 x3 0 −1 1
−1 2 1 5 2
4 3 5 x4 4 3 5
11 6 −4 8 x1 8
6 12 0 24 x2 12
= .
−4 0 2 2 x3 −1
8 24 2 50 x4 23
This operation has mapped both sides of the equation into the same space, namely, the column space
of AT , which is also the row space of A. Since the rank of A is r = 2, the dimension of the row space
is also two, and now the vectors on both sides of the equation have been mapped into the same plane.
Again using row operations to perform Gaussian elimination gives rise to
1 0 −1/2 −1/2 x1 1/4
0 1 1/4 9/4 x2 7/8
= .
0 0 0 0 x3 0
0 0 0 0 x4 0
This equation suggests that here x3 and x4 are arbitrary, so we set x3 = s, x4 = t and, treating s and
t as known quantities, reduce the system to the following
1 0 x1 1/4 + s/2 + t/2
= ,
0 1 x2 7/8 − s/4 − 9t/4
so
x1 1/4 1/2 1/2
x2 7/8 −1/4 −9/4
= + s + t .
x3 0 1 0
x4 0 0 1
The vectors which are multiplied by s and t are in the right null space of A. The vector (1/4, 7/8, 0, 0)T
is not entirely in the row space of A; it has components in both the row space and right null space. We
can, thus, decompose this vector into a linear combination of row space vectors and right null space
vectors using the procedure in the previous section, solving the following equation for the coefficients
a1 , . . . , a4 , which are the coefficients of the row and right null space vectors:
1/4 1 3 1/2 1/2 a1
7/8 2 2 −1/4 −9/4 a2
= .
0 0 −1 1 0 a3
0 4 3 0 1 a4
Solving, we get
a1 −3/244
a2 29/244
= .
a3 29/244
a4 −75/244
So we can recast the solution as
x1 1 3 1/2 1/2
x2 3 2 29 2 29 −1/4 75 −9/4
=− + + s + + t− .
x3 244 0 244 −1 244 1 244 0
x4 4 3 0 1
| {z } | {z }
row space right null space
This choice of x guarantees that we minimize ||A · x − b||2 , which in this case is 1.22474. So there are
no vectors x which satisfy the original equation A · x = b, but there are a doubly infinite number of
vectors x which can minimize the least squares error.
We can choose special values of s and t such that we minimize ||x||2 while maintaining ||A · x − b||2
at its global minimum. This is done simply by forcing the magnitude of the right null space vectors to
zero, so we choose s = −29/244, t = 75/244, giving
x1 1 3 21/61
x2 3 2 29 2 13/61
=− + = .
x3 244 0 244 −1 −29/244
x4 4 3 75/244
| {z }
row space
det Ai
xi = (8.13)
det A
where Ai is the matrix obtained by replacing the i-th column of A by y. While generally
valid, Cramer’s rule is most useful for low dimension systems. For large systems, Gaussian
elimination, is a more efficient technique.
Example 8.11
For A: R2 → R2 , Solve for x in A · x = b:
1 2 x1 4
=
3 2 x2 5
By Cramer’s rule
4 2
5 2 −2 1
x1 = = =
1 2
−4 2
3 2
1 4
3 5 −7 7
x2 = = =
1 2
−4 4
3 2
So 1
x= 2
7
4
We get the same result by Gaussian elimination. Subtracting three times the first row from the second
yields
1 2 x1 4
=
0 −4 x2 −7
Thus x2 = 74 . Back substitution into the first equation then gives x1 = 12 .
Example 8.12
With A : R2 → R2 , find the most general x which best satisfies A · x = b for
1 2 x1 2
= .
3 6 x2 0
Obviously, there is no unique solution to this system since the determinant of the coefficient matrix is
zero. The rank of A is 1, so in actuality, A maps vectors from R2 into a one dimensional subspace, R1 .
For a general b, which does not lie in the one dimensional subspace, we can find the best solution x by
first multiplying both sides by AT :
1 3 1 2 x1 1 3 2
= .
2 6 3 6 x2 2 6 0
10 20 x1 2
= .
20 40 x2 4
This operation maps both sides of the equation into the column space of AT , which is the row space
of A, which has dimension 1. Since the vectors are now in the same space, a solution can be found.
Using row reductions to perform Gaussian elimination, we get
1 2 x1 1/5
= .
0 0 x2 0
The vector which t multiplies, (−2, 1)T , is in the right null space of A. We can recast the vector
(1/5, 0)T in terms of a linear combination of the row space vector (1, 2)T and the right null space vector
to get the final form of the solution:
x1 1 1 2 −2
= + t− .
x2 25 2 25 1
| {z } | {z }
row space right null space
This choice of x guarantees that the least squares error ||A · x − b||2 is minimized. In this case the least
squares error is 1.89737. The vector x with the smallest norm that minimizes ||A · x − b||2 is found by
setting the magnitude of the right null space contribution to zero, so we can take t = 2/25 giving
x1 1 1
= .
x2 25 2
| {z }
row space
that it is the right set that is being discussed. Though it does not often arise, there are
occasions when one requires the left eigenvectors which arise from eT · A = eT · Iλ. If the
matrix A is self-adjoint, it can be shown that it has the same left and right eigenvectors.
If A is not self-adjoint, it has different left and right eigenvectors. The eigenvalues are the
same for both left and right eigenvectors of the same operator, whether or not the system is
self-adjoint.
Second, the polynomial equation that arises in the eigenvalue problem is the characteristic
equation of the matrix.
Theorem
A matrix satisfies its own characteristic equation (Cayley-Hamilton1 theorem).
If a matrix is triangular, then the eigenvalues are the diagonal terms. Eigenvalues of
A2 are the square of the eigenvalues of A. Every eigenvector of A is also an eigenvector
of A · A = A2 . The spectral radius of a matrix is the largest of the absolute values of the
eigenvalues.
The trace of a matrix is the sum of the terms on the leading diagonal.
Theorem
The trace of a n × n matrix is the sum of its n eigenvalues.
Theorem
The product of the n eigenvalues is the determinant of the matrix.
Example 8.13
Demonstrate the above theorems for
0 1 −2
A = 2 1 0 .
4 −2 5
The characteristic equation is
λ3 − 6λ2 + 11λ − 6 = 0.
The Cayley-Hamilton theorem is easily verified by direct substitution:
A3 − 6A2 + 11A − 6I = 0,
0 1 −2 0 1 −2 0 1 −2 0 1 −2 0 1 −2
2 1 0 2 1 0 2 1 0 − 62 1 0 2 1 0
4 −2 5 4 −2 5 4 −2 5 4 −2 5 4 −2 5
0 1 −2 1 0 0 0 0 0
+11 2 1 0 − 60 1 0 = 0 0 0,
4 −2 5 0 0 1 0 0 0
−30 19 −38 36 −30 60 0 11 −22 −6 0 0
−10 13 −24 + −12 −18 24 + 22 11 0 + 0 −6 0
52 −26 53 −96 48 −102 44 −22 55 0 0 −6
1
after Arthur Cayley, 1821-1895, English mathematician, and William Rowan Hamilton, 1805-1865,
Anglo-Irish mathematician.
0 0 0
= 0 0 0.
0 0 0
Considering the traditional right eigenvalue problem, A·e = λI·e, it is easily shown that the eigenvalues
and (right) eigenvectors for this system are
0
λ1 = 1, e(1) = 2 ,
1
1
2
λ2 = 2, e(2) = 1 ,
0
−1
λ3 = 3, e(3) = −1 .
1
One notes that while the eigenvectors do form a basis in R3 , that they are not orthogonal; this is a
consequence of the matrix not being self-adjoint (or more specifically asymmetric). The spectral radius
of A is 3. Now
0 1 −2 0 1 −2 −6 5 −10
A2 = A · A = 2 1 0 2 1 0 = 2 3 −4 .
4 −2 5 4 −2 5 16 −8 17
It is easily shown that the eigenvalues for A2 are 1, 4, 9, precisely the squares of the eigenvalues of A.
The trace is
trA = 0 + 1 + 5 = 6.
Note this is the equal to the sum of the eigenvalues
3
X
λi = 1 + 2 + 3 = 6.
i=1
In this case As has real eigenvalues, both positive and negative, λ1 = 5.32, λ2 = −1.39, λ3 = 2.07.
Because of the presence of a negative eigenvalue in the symmetric part of A, we can conclude that both
A and As are not positive definite.
We also note that for real-valued problems x ∈ RN , A ∈ RN ×N , the antisymmetric part of a matrix
can never be positive definite by the following argument. We can say xT · A · x = xT · (As + Aa ) · x.
Then one has xT · Aa · x = 0 for all x because the tensor inner product of the real antisymmetric Aa
with the symmetric xT and x is identically zero. So to test the positive definiteness of a real A, it
suffices to consider the positive definiteness of its symmetric part: xT · As · x ≥ 0.
Example 8.14
If
1+i
3 − 2i
x= ,
2
−3i
find ||x||2 .
v
u
u 1+i
√ u 3 − 2i √ √
||x||2 = + xH · x = +u
t(1 − i, 3 + 2i, 2, +3i) 2 = + 2 + 13 + 4 + 9 = 2 7.
−3i
Example 8.15
If
1+i
x = −2 + 3i ,
2−i
3
y = 4 − 2i ,
3 + 3i
find <x, y>.
<x, y> = xH · y,
3
= (1 − i, −2 − 3i, 2 + i) 4 − 2i ,
3 + 3i
= (3 − 3i) + (−14 − 8i) + (3 + 9i),
= −8 − 2i.
AH = ĀT .
As the Hermitian transpose is the adjoint operator corresponding to a given complex matrix,
we can apply an earlier proved theorem for linear operators to deduce that the eigenvalues
of a complex matrix are the complex conjugates of the Hermitian transpose of that matrix.
The Hermitian transpose is distinguished from a matrix which is Hermitian as follows. A
Hermitian matrix is one which is equal to its conjugate transpose. So a matrix which equals
its Hermitian transpose is Hermitian. A matrix which does not equal its Hermitian transpose
is non-Hermitian. A skew-Hermitian matrix is the negative of its Hermitian transpose. A
Hermitian matrix is self-adjoint.
Properties:
• xH · A · x is real if A is Hermitian.
Note the diagonal elements of a Hermitian matrix must be real as they must be unchanged
by conjugation.
Example 8.16
Consider A · x = b, where A : C3 → C3 with A the Hermitian matrix and x the complex vector:
1 2−i 3 3 + 2i
A = 2 + i −3 2i , x = −1 .
3 −2i 4 2−i
First, we have
1 2−i 3 3 + 2i 7
b = A · x = 2 + i −3 2i −1 = 9 + 11i .
3 −2i 4 2−i 17 + 4i
Now, demonstrate that the properties of Hermitian matrices hold for this case. First
1 2−i 3 3 + 2i
xH · A · x = (3 − 2i, −1, 2 + i) 2 + i −3 2i −1 = 42 ∈ R1 .
3 −2i 4 2−i
Check for orthogonality between two of the eigenvectors, e.g e(1) , e(2) :
Example 8.17
Find the orthogonal matrix corresponding to
2 1
A= .
1 2
1 1 1 1
The normalized eigenvectors are 2
√ and 2
√ .
−1 1
The orthogonal matrix is
1 1 1
Q= √ .
2 −1 1
Note that Q is not symmetric.
Example 8.18
Consider the unitary matrix !
1+i
√ 1−2i
√
3 15
U= .
√1 1+3i
√
3 15
The column vectors are easily seen to be normal. They are also orthogonal:
1−2i !
1−i 1 √
15
√ ,√ 1+3i = 0 + 0i.
3 3 √
15
The matrix itself is not Hermitian. Still, its Hermitian transpose exists:
1−i 1
!
√ √
UH = 3
1+2i
3
1−3i .
√ √
15 15
Here κ is the wavenumber, and is the reciprocal of the wavelength. The FT has a discrete
analog. The connection between the two is often not transparent in the literature. With some
effort a connection can be made at the expense of diverging from one school’s notation to
the other’s. Here, we will be satisfied with a form which demonstrates the analogs between
the continuous and discrete transform, but will not be completely linked. To make the
connection, one can construct a discrete approximation to the integral of the FT, and with
some effort, arrive at an equivalent result.
For the DFT, consider a function y(x), x ∈ [xmin , xmax ], x ∈ R1 , y ∈ R1 . Now discretize
the domain into N uniformly distributed points so that every xj is mapped to a yj for
j = 0, . . . , N −1. Here we comply with the tradtional, yet idiosyncratic, limits on j which are
found in many texts on DFT. This offsets standard vector and matrix numbering schemes by
one, and so care must be exercised in implementing these algorithms with common software.
We seek a discrete analog of the continuous Fourier transformation of the form
N −1
1 X N −1 xj − xmin
yj = √ ck exp (2πi)k , j = 0, . . . , N − 1. (8.17)
N k=0 N xmax − xmin
Here k plays the role of κ, and ck plays the role of Y (κ). For uniformly spaced xj , one has
xj − xmin
j = (N − 1) , (8.18)
xmax − xmin
FH · F = I. (8.24)
Since F is unitary, it is immediately known that FH = F−1 , that ||F||2 = 1, that the
eigenvalues of F have magnitude of unity, and that the column vectors of F are orthonormal.
√
Note that F is not Hermitian. Also note that many texts omit the factor 1/ N in the
definition of F; this is not a major problem, but does render F to be non-unitary.
Now given a vector y = yj , j = 0, . . . , N −1, the DFT is defined as the following mapping
c = FH · y. (8.25)
F·c = F · FH · y, (8.26)
F·c = F · F−1 · y, (8.27)
F·c = I · y, (8.28)
y = F · c. (8.29)
Now ||F||2 = 1 because it is an orthogonal matrix, and because all its eigenvalues have unit
magnitude, it in fact does not change the norm of c, so we get the equality, which induces a
Parseval’s equation
||y||2 = ||c||2. (8.32)
Example 8.19
Consider a five term DFT of the function
Take then for N = 5, a set of uniformly distributed points in the domain and their image in the range:
x0 = 0, x1 = 1, x2 = 2, x3 = 3, x4 = 4, (8.34)
y0 = 0, y1 = 1, y2 = 4, y3 = 9, y4 = 16. (8.35)
z (0) = w0 = 1, (8.37)
z (1) = w1 = 0.3090 + 0.9511i, (8.38)
z (2) = w2 = −0.8090 + 0.5878i, (8.39)
z (3) = w3 = −0.8090 − 0.5878i, (8.40)
z (4) = w4 = 0.3090 − 0.9511i, (8.41)
(8.42)
Now c = FH · y, so
c0 1 1 1 1 1 0
c1 1 0.3090 − 0.9511i −0.8090 − 0.5878i −0.8090 + 0.5878i 0.3090 + 0.9511i 1
1
c2 = √ 1 −0.8090 − 0.5878i 0.3090 + 0.9511i 0.3090 − 0.9511i −0.8090 + 0.5878i 4
c3 5 1 −0.8090 + 0.5878i 0.3090 − 0.9511i 0.3090 + 0.9511i −0.8090 − 0.5878i
9
c4 1 0.3090 + 0.9511i −0.8090 + 0.5878i −0.8090 − 0.5878i 0.3090 − 0.9511i 16
13.4164
−2.3541 + 7.6942i
= −4.3541 + 1.8164i (8.45)
−4.3541 − 1.8164i
−2.3541 − 7.6942i
Now one is often interested in the magnitude of the components of c, which gives a measure of the
so-called energy associated with each Fourier mode. So one calculates a vector of the magnitude of
each component as
√
√c0 c0 |c0 | 13.4164
√c1 c1 |c1 | 8.0463
√c2 c2 = |c2 | = 4.7178 (8.46)
√ c 3 c 3 |c 3 | 4.7178
c4 c4 |c4 | 8.0463
Now due to a phenomena known as aliasing, explained in detail in standard texts, the values of ck
which have the most significance are the first half ck , k = 0, . . . , N/2.
Here √
||y||2 = ||c||2 = 354 = 18.8149. (8.47)
Note that by construction
1
y0 = √ (c0 + c1 + c2 + c3 + c4 ) , (8.48)
5
1
y1 = √ c0 + c1 e2πi/5 + c2 e4πi/5 + c3 e6πi/5 + c4 e8πi/5 , (8.49)
5
1
y2 = √ c0 + c1 e4πi/5 + c2 e8πi/5 + c3 e12πi/5 + c4 e16πi/5 , (8.50)
5
1
y3 = √ c0 + c1 e6πi/5 + c2 e12πi/5 + c3 e18πi/5 + c4 e24πi/5 , (8.51)
5
1
y4 = √ c0 + c1 e8πi/5 + c2 e16πi/5 + c3 e24πi/5 + c4 e32πi/5 . (8.52)
5
In general, it is seen that yj can be described by
N −1
1 X kj
yj = √ ck exp (2πi) , j = 0, . . . , N − 1. (8.53)
N k=0 N
Realizing now that for a uniform discretization, such as done here, that
xmax − xmin
∆x = , (8.54)
N −1
and that
xj = j∆x + xmin , j = 0, . . . , N − 1, (8.55)
one has
xmax − xmin
xj = j + xmin , j = 0, . . . , N − 1. (8.56)
N −1
Solving for j, one gets
xj − xmin
j = (N − 1) , (8.57)
xmax − xmin
so that yj can be expressed as a Fourier-type expansion in terms of xj as
N
1 X N −1 xj − xmin
yj = √ ck exp (2πi)k , j = 0, . . . , N − 1. (8.58)
N k=1 N xmax − xmin
N −1
κk = k . (8.59)
N
And as N → ∞, one has
κk ∼ k. (8.60)
Example 8.20
The real power of the DFT is seen in its ability to select ampitudes of modes of signals at certain
frequencies. Consider the signal for x ∈ [0, 3]
2x 10x 100x
y(x) = 10 sin (2π) + 2 sin (2π) + sin (2π) . (8.61)
3 3 3
Rescaling the domain so as to take x ∈ [0, 3] into x̃ ∈ [0, 1] via the transformation x̃ = x/3, one has
To capture the high wavenumber components of the signal, one must have a suffiently large value of N .
Note in the transformed domain that the smallest wavelength is λ = 1/100 = 0.01. So for a domain
length of unity, one needs at least N = 100 sampling points. In fact, let us choose to take more points,
N = 523. There is no problem in choosing an unusual number of points for this so-called slow Fourier
transform. If an FFT were attempted, one would have to choose integral powers of 2 as the number of
points.
A plot of the function y(x) and two versions of its DFT, |ck | vs. k, is given in in Figure 8.3 Note
that |ck | has its peaks at k = 2, k = 10, and k = 100, equal to the wavenumbers of the generating sine
functions, κ1 = 2, κ2 = 10, and κ3 = 100. To avoid the confusing, and non-physical, aliasing effect,
only half the |ck | values have been plotted the first DFT of Fig. 8.3. The second DFT here plots all
values of |ck | and thus exhibits aliasing for large k.
Original Signal
15
10
0
y
−5
−10
−15
0 1 2 3
x
100 100
80 80
|c |
|ck|
k
60 60
40 40
20 20
0 0
0 1 2 3 −200 0 200 400 600
10 10 10 10
k k
Figure 8.3: Plot of a three term sinusoid y(x) and its discrete Fourier transform for N = 523
points. The first DFT is plotted from k = 0, . . . , N/2 and thus represents the original signal
well. The second DFT is plotted from k = 0, . . . , N − 1 and exhibits aliasing effects at high
k.
Original Signal
3
0
y
−1
−2
−3
0 0.2 0.4 0.6 0.8 1
x
12 12
10 10
8 8
|ck|
|c |
k
6 6
4 4
2 2
0 0
0 1 2 3
10 10 10 10 −200 0 200 400 600 800
k k
Figure 8.4: Plot of a two term sinusoid accompanied by random noise y(x) and its discrete
Fourier transform for N = 607 points. The first DFT is plotted from k = 0, . . . , N/2 and
thus represents the original signal well. The second DFT is plotted from k = 0, . . . , N − 1
and exhibits aliasing effects at high k.
Example 8.21
Now take the DFT of a signal which is corrupted by so-called white, or random, noise. The signal
here is given in x ∈ [0, 1] by
Here frand [−1, 1](x) returns a random number between −1 and 1 for any value of x. A plot of the
function y(x) and two versions of its 607 point DFT, |ck | vs. k, is given in in Figure 8.4 In the raw data
plotted in Fig. 8.4, it is difficult to distinguish the signal from the random noise. But on examination
of the accompanying DFT plot, it is clear that there are unambiguous components of the signal which
peak at k = 10 and k = 100, which indicates there is a strong component of the signal with κ = 10 and
κ = 100. Once again, to avoid the confusing, and non-physical, aliasing effect, only half the |ck | values
have been plotted in the first DFT of Fig. 8.4. The second DFT gives all values of |ck | and exhibits
aliasing.
8.5.1 L · D · U decomposition
Probably the most important technique in solving linear systems of algebraic equations of
the form A · x = b, uses the decomposition
A = P−1 · L · D · U, (8.64)
P−1 · L · D · U · x = b.
P−1 · L · D · c = b
so −1
c = P−1 · L · D ·b
The triangular form of L · D renders the inversion of (P−1 · L · D) to be much more
computationally efficient than inversion of an arbitrary square matrix.
2
If A is not square, there is an equivalent decomposition, known as row echelon form, to be discussed
later in this chapter.
U·x=c
so
x = U−1 · c
Again since U is triangular, the inversion is computationally efficient.
Example 8.22
Find the L · D · U decomposition of the matrix below:
−5 4 9
A = −22 14 18 .
16 −8 −6
The process is essentially a series of row operations, which is the essence of Gaussian elimination. First
we operate to transform the −22 and 16 in the first column into zeroes. Crucial in this step is the
necessity of the term in the 1,1 slot, known as the pivot, to be non-zero. If it is zero, a row exchange
will be necessary, mandating a permutation matrix which is not the identity matrix. In this case there
are no such problems. We multiply the first row by 22 5 and subtract from the second row, then multiply
the first row by − 16
5 and subtract from the third row. The factors 22 16
5 and − 5 will go in the 2,1 and
3,1 slots of the matrix L. The diagonal of L always is filled with ones. This row operation yields
−5 4 9 1 0 0 −5 4 9
A = −22 14 18 = 22/5 1 0 0 −18/5 −108/5 .
16 −8 −6 −16/5 0 1 0 24/5 114/5
Now multiplying the new second row by − 34 , subtracting this from the third row, and depositing the
factor − 34 into 3,2 slot of the matrix L, we get
−5 4 9 1 0 0 −5 4 9
A = −22 14 18 = 22/5 1 0 0 −18/5 −108/5 .
16 −8 −6 −16/5 −4/3 1 0 0 −6
| {z }| {z }
L U
The form above is often described as the L · U decomposition of A. We can force the diagonal terms
of the upper triangular matrix to unity by extracting a diagonal matrix D to form the L · D · U
decomposition:
−5 4 9 1 0 0 −5 0 0 1 −4/5 −9/5
A = −22 14 18 = 22/5 1 0 0 −18/5 0 0 1 6 .
16 −8 −6 −16/5 −4/3 1 0 0 −6 0 0 1
| {z }| {z }| {z }
L D U
Note that D does not contain the eigenvalues of A. Also since there were no row exchanges necessary
P = P−1 = I, and it has not been included.
Example 8.23
Find the L · D · U decomposition of the matrix A:
0 1 2
A = 1 1 1.
1 0 0
Performing Gaussian elimination by subtracting 1 times the first row from the second and depositing
the 1 in the 2,1 slot of L, we get
0 0 1 0 1 2 1 0 0 1 0 0
P ·A = L· U → 0 1 01 1 1 = 1 1 00 1 1.
1 0 0 1 0 0 0 0 1 0 1 2
Now subtracting 1 times the second row, and depositing the 1 in the 3,2 slot of L
0 0 1 0 1 2 1 0 0 1 0 0
P · A = L · U → 0 1 01 1 1 = 1 1 00 1 1.
1 0 0 1 0 0 0 1 1 0 0 1
Now U already has ones on the diagonal, so the diagonal matrix D is simply the identity matrix. Using
this and inverting P, which is P itself(!), we get the final decomposition
0 1 2 0 0 1 1 0 0 1 0 0 1 0 0
A = P−1 · L · D · U → 1 1 1 = 0 1 0 1 1 0 0 1 0 0 1 1 .
1 0 0 1 0 0 0 1 1 0 0 1 0 0 1
| {z }| {z }| {z }| {z }
P−1 L D U
It can also be shown that if A is symmetric, that the decomposition can be written as
A = P−1 · L · D · LT .
Example 8.24
Consider the non-square matrix studied earlier,
1 −3 2
A= .
2 0 3
We take 2 times the first row and subtract the result from the second row. The scalar 2 is deposited in
the 2,1 slot in the L matrix. So
1 −3 2 1 0 1 −3 2
A= = .
2 0 3 2 1 0 6 −1
| {z } | {z }
L U
Again, the above is also known as an L · U decomposition, and is often as useful as the L · D · U
decomposition. There is no row exchange so the permutation matrix and its inverse are the identity
matrix. We extract a 1 and 6 to form the diagonal matrix D, so the final form is
1 0 1 0 1 0 1 −3 2
A = P−1 · L · D · U = .
0 1 2 1 0 6 0 1 − 16
| {z } | {z } | {z } | {z }
P−1 L D U
Row echelon form is an especially useful form for underconstrained systems as illustrated
in the following example.
Example 8.25
Consider solutions for the unknown x in the equation A · x = b where A is known A : R5 → R3 ,
and b is left general, but considered to be known:
x1
2 1 −1 1 2 x2 b1
4 2 −2 1 0 x3 = b2 .
−2 −1 1 −2 −6 x4 b3
x5
We perform Gaussian elimination row operations on the second and third rows to get zeros in the first
column:
x1
2 1 −1 1 2 x2 b1
0 0 0 −1 −4
x3 = −2b1 + b2 .
0 0 0 −1 −4 x4 b1 + b3
x5
The next round of Gaussian elimination works on the third row and yields
x1
2 1 −1 1 2 x2 b1
0 0 0 −1 −4
x3 = −2b1 + b2 .
0 0 0 0 0 x4 3b1 − b2 + b3
x5
0 = 3b1 − b2 + b3 .
This is the equation of a plane in R3 . Thus arbitrary b ∈ R3 will not satisfy the original equation.
Said another way, the operator A maps arbitrary five dimensional vectors x into a two-dimensional
subspace of a three dimensional vector space. The rank of A is 2. Thus the dimension of both the row
space and the column space is 2; the dimension of the right null space is 3, and the dimension of the
left null space is 1.
We also note there are two non-trivial equations remaining. The first non-zero elements from the
left of each row are known as the pivots. The number of pivots is equal to the rank of the matrix.
Variables which correspond to each pivot are known as basic variables. Variables with no pivot are
known as free variables. Here the basic variables are x1 and x4 , while the free variables are x2 , x3 , and
x5 .
Now enforcing the constraint 3b1 − b2 + b3 = 0, without which there will be no solution, we can
set each free variable to an arbitrary value, and then solve the resulting square system. Take x2 = r,
x3 = s, x5 = t, where here r, s, and t are arbitrary real scalar constants. So
x1
2 1 −1 1 2 r b1
0 0 0 −1 −4 s = −2b1 + b2 ,
0 0 0 0 0 x4 0
t
which gives
2 1 x1 b1 − r + s − 2t
= ,
0 −1 x4 −2b1 + b2 + 4t
which yields
x4 = 2b1 − b2 − 4t,
1
x1 = (−b1 + b2 − r + s + 2t).
2
Thus
1 1 1 1
x1 2 (−b1 + b2 − r + s + 2t) 2 (−b1 + b2 ) −2 2 1
x2 r 0 1 0 0
x = x3 = s = 0 + r 0 + s 1 + t 0
x4 2b1 − b2 − 4t 2b1 − b2 0 0 −4
x5 t 0 0 0 1
r, s, t ∈ R1 .
The coefficients r, s, and t multiply the three right null space vectors. These in combination with two
independent row space vectors, form a basis for any vector x. Thus, we can again cast the solution as a
particular solution which is a unique combination of independent row space vectors and a non-unique
combination of the right null space vectors (the homogeneous solution):
1 1
x1 2 4 −2 2 1
x2 25b − 13b 1 −13b + 11b 2 1 0 0
1 2 1 2
x = x3 = −1 + −2 + r̂ 0 + ŝ 1 + t̂ 0 .
106 106
x4 1 1 0 0 −4
x5 2 0 0 0 1
| {z } | {z }
row space right null space
The L · D · U decomposition is
1
2 1 −1 1 2 1 0 0 2 0 0 1 2 − 21 1
2 1
4 2 −2 1 0 = 2 1 00 −1 0 0 0 0 1 4.
−2 −1 1 −2 −6 −1 1 1 0 0 0 0 0 0 0 0
| {z } | {z }| {z }| {z }
A L D U
There were no row exchanges, so in effect the permutation matrix P is the identity matrix, and there
is no need to include it.
Lastly, we note that a more robust alternative to the method shown here would be to first apply
the AT operator to both sides of the equation so to map both sides into the column space of A. Then
there would be no need to restrict b so that it lies in the column space. Our results are then interpreted
as giving us only a projection of x. Taking AT · A · x = AT · b and then casting the result into row
echelon form gives
1 1/2 −1/2 0 −1 x1 (1/22)(b1 + 7b2 + 4b3 )
0 0 0 1 4 x2 (1/11)(b1 − 4b2 − 7b3 )
0 0 0 0 0 x3 = 0 .
0 0 0 0 0 x4 0
0 0 0 0 0 x5 0
This suggests we take x2 = r, x3 = s, and x5 = t and solve so to get
x1 (1/22)(b1 + 7b2 + 4b3 ) −1/2 1/2 1
x2 0 1 0 0
x3 = 0 + r 0 + s 1 + t 0 .
x4 (1/11)(b1 − 4b2 − 7b3 ) 0 0 −4
x5 0 0 0 1
We could go on to cast this in terms of combinations of row vectors and right null space vectors, but
will not do so here. It is reiterated the this result is valid for arbitrary b, but that it only represents a
solution which minimizes the error in ||A · x − b||2 .
8.5.3 Q · R decomposition
The Q · R decomposition allows us to formulate and matrix as the product of an orthogonal
(unitary if complex) matrix Q and an upper triangular matrix R, of the same dimension as
A. That is we seek Q and R such that
A=Q·R (8.65)
The matrix A can be square or rectangular. See Strang for details of the algorithm.
Example 8.26
The Q · R decomposition of the matrix we diagonalized in a previous example is as follows:
−5 4 9 −0.1808 −0.4982 0.8480 27.6586 −16.4867 −19.4153
A = −22 14 18 = Q · R = −0.7954 −0.4331 −0.4240 0 −2.0465 −7.7722 .
16 −8 −6 0.5785 −0.7512 −0.3180 0 0 1.9080
| {z } | {z }| {z }
A Q R
Noting that ||Q||2 = 1, we deduce that ||R||2 = ||A||2. Also recalling how matrices can
be thought of as transformations, we see how to think of A as a pure rotation (Q) followed
by stretching (R).
Example 8.27
Find the Q · R decomposition for our non-square matrix
1 −3 2
A= .
2 0 3
The decomposition is
−0.4472 −0.8944 −2.2361 1.3416 −3.577
A= .
−0.8944 0.4472 0 2.6833 −0.4472
| {z }| {z }
Q R
The Q·R decomposition can be shown to be closely related to the Gram-Schmidt orthog-
onalization process. It is also useful in increasing the efficiency of estimating x for A · x ≃ b
when the system is overconstrained; that is b is not in the column space of A, R(A). If we,
−1
A · x = Q · R · RT · R · RT · QT · b.
−1
When rectangular R has no zeros on its diagonal, R · RT · R · RT has all zeroes, except
for r ones on the diagonal, where r is the rank of R. This makes solution of overconstrained
problems particularly simple.
8.5.4 Diagonalization
Casting a matrix into a form in which all (or sometimes most) of its off-diagonal elements
have zero value has its most important application in solving systems of differential equations
but also in other scenarios. For many cases, we can decompose a square matrix A into the
form
A = S · Λ · S−1 , (8.66)
where S is non-singular matrix and Λ is a diagonal matrix. To diagonalize a square matrix
A, we must find S, a diagonalizing matrix, such that S−1 · A · S is diagonal. Not all matrices
are diagonalizable.
Theorem
A matrix with distinct eigenvalues can be diagonalized, but the diagonalizing matrix is
not unique.
Definition: The algebraic multiplicity of an eigenvalue is number of times it occurs. The
geometric multiplicity of an eigenvalue is the number of eigenvectors it has.
Theorem
Nonzero eigenvectors corresponding to different eigenvalues are linearly independent.
Theorem
If A is an n × n matrix with n linearly independent right eigenvectors {e1 , e2 , · · · , en }
corresponding to eigenvalues {λ1 , λ2 , · · · , λn } (not necessarily distinct), then the n×n matrix
S whose columns are populated by the eigenvectors of A
S = {e1 }, {e2 }, · · · , {en } (8.67)
makes
S−1 · A · S = Λ, (8.68)
where
λ1 0 · · · 0
0 λ2 · · · 0
Λ= .. .. .. .. (8.69)
. . . .
0 0 · · · λn
is a diagonal matrix of eigenvalues. The matrices A and Λ are similar.
Let’s see if this recipe works when we fill the columns of S with the eigenvectors.
S−1 · A · S = Λ, (8.70)
A · S = S · Λ, (8.71)
(1) (n) (1) (n)
a11 · · · a1n e1 · · · e1 e1 · · · e1 λ1 ··· 0
... ..
.
.. ..
. .
..
.
.. = ..
. .
..
.
.. ..
. .
.. .
. .. ,(8.72)
(1) (n) (1) (n)
an1 · · · ann e · · · en en · · · en 0 · · · λn
| {z }| n {z } | {z }| {z }
=A =S =S =Λ
(1) (n)
λ1 e1 ··· λn e1
= .. .. .. , (8.73)
. . .
(1) (n)
λ1 en · · · λn en
| {z }
=S·Λ
(1) (1)
A·e = λ1 e , (8.74)
(2)
A·e = λ2 e(2) , (8.75)
.. ..
. .
(n)
A·e = λn e(n) . (8.76)
A · S · S−1 = S · Λ · S−1 ,
A = S · Λ · S−1 .
Example 8.28
Diagonalize the matrix considered in a previous example:
−5 4 9
A = −22 14 18 ,
16 −8 −6
and check.
Then
−1 1 2
S = ( e1 e2 e3 ) = −2 2 1 .
1 0 2
The inverse is
−4 2 3
1
S−1 = −5 4 3 .
3
2 −1 0
Thus
6 3 12
A · S = 12 6 6 ,
−6 0 12
and
−6 0 0
Λ = S−1 · A · S = 0 3 0 .
0 0 6
Let us also note the complementary decomposition of A:
−1 1 2 −6 0 0 −4 2 3 −5 4 9
1
A = S · Λ · S−1 = −2 2 1 0 3 0 −5 4 3 = −22 14 18 .
3
1 0 2 0 0 6 2 −1 0 16 −8 −6
| {z }| {z }| {z } | {z }
S Λ S−1 A
Note that because the matrix is not symmetric, the eigenvectors are not orthogonal, e.g. eT1 · e2 = −5.
Note that if A is symmetric (Hermitian), then its eigenvectors must be orthogonal; thus,
it is possible to normalize the eigenvectors so that the matrix S is in fact orthogonal (unitary
if complex). Thus for symmetric A we have
A = Q · Λ · Q−1 .
Note also that with A · S = S · Λ, the column vectors of S (which are the right eigenvectors
of A) form a basis in Cn .
Consider now the right eigensystem of the adjoint of A, denoted by A∗ :
A∗ · P = P · Λ ∗ (8.77)
where Λ∗ is the diagonal matrix containing the eigenvalues of A∗ , and P is the matrix whose
columns are populated by the (right) eigenvectors of A∗ . Now we know from an earlier
proof that the eigenvalues of the adjoint are the complex conjugates of those of the original
operator, thus Λ∗ = ΛH . Also the adjoint operator for matrices is the Hermitian transpose.
So, we find that
AH · P = P · Λ H (8.78)
Taking the Hermitian transpose of both sides, we recover
PH · A = Λ · PH . (8.79)
So we see clearly that the left eigenvectors of a linear operator are the right eigenvectors of
the adjoint of that operator.
It is also possible to show that, remarkably, when we take the inner product of the matrix
of right eigenvectors of the operator with the matrix of right eigenvectors of its adjoint, that
we obtain a diagonal matrix, which we denote as D:
SH · P = D. (8.80)
Equivalently, this states that the inner product of the left eigenvector matrix with the right
eigenvector matrix is diagonal. Let us see how this comes about. Let si be a right eigenvector
of A with eigenvalue λi and pj be a left eigenvector of A with eigenvalue λj . Then
A · si = λi si , (8.81)
and
pH H
j · A = λj pj . (8.82)
If we premultiply the first eigen-relation, Eq. (8.81), by pH
j , we obtain
pH · A ·si = pH
j · (λi si ) . (8.83)
| j {z }
=λj pH
j
Substituting from the second eigen-relation, Eq. (8.82) and rearranging, Eq. (8.83) becomes
λj pH H
j · si = λi pj · si . (8.84)
Rearranging
(λj − λi ) pH
j · si = 0. (8.85)
pH
j · si = 0, (8.86)
sH
i · pj = 0. (8.87)
SH · P̂ = I. (8.88)
Here P̂ denotes the matrix in which each eigenvector (column) of the original P has been
scaled such that the above identity is achieved. Hence P̂ is seen to give the set of reciprocal
basis vectors for the basis defined by S:
SR = P̂. (8.89)
It is also easy to see then that the inverse of the matrix S is given by
Example 8.29
For a matrix A considered in an earlier example, consider the basis formed by its matrix of eigen-
vectors S, and use the properly scaled matrix of eigenvectors of A∗ = AH to determine the reciprocal
basis SR .
We will take
−5 4 9
A = −22 14 18 .
16 −8 −6
As found before, the eigenvalue- (right) eigenvector pairs are
−1
λ1 = −6, e1R = −2 ,
1
1
λ2 = 3, e2R = 2 ,
0
2
λ3 = 6, e3R = 1 .
2
The adjoint of A is
−5 −22 16
AH = 4 14 −8 .
9 18 −6
The eigenvalues-(right) eigenvectors of AH , which are the left eigenvectors of A, are found to be
−4
λ1 = −6, e1L = 2 ,
3
−5
λ2 = 3, e2L = 4 ,
3
−2
λ3 = 6, e3L = 1 .
0
So the matrix of right eigenvectors of the adjoint, which contains the left eigenvectors of the original
matrix, is
−4 −5 −2
P= 2 4 1 .
3 3 0
We indeed find that the inner product of S and P is a diagonal matrix D:
−1 −2 1 −4 −5 −2 3 0 0
SH · P = 1 2 0 · 2 4 1 = 0 3 0 .
2 1 2 3 3 0 0 0 −3
Using our knowledge of D, we individually scale each column of P to form the desired reciprocal basis
−4/3 −5/3 2/3
P̂ = 2/3 4/3 −1/3 = SR .
1 1 0
Then we see that the inner product of S and the reciprocal basis P̂ = SR is indeed the identity matrix:
−1 −2 1 −4/3 −5/3 2/3 1 0 0
SH · P̂ = 1 2 0 · 2/3 4/3 −1/3 = 0 1 0 .
2 1 2 1 1 0 0 0 1
Example 8.30
Find the Jordan canonical form of
4 1 3
A= 0 4 1
0 0 4
3
Marie Ennemond Camille Jordan, 1838-1922, French mathematician.
A = Q · R · QT . (8.97)
Here Q is an orthogonal (unitary if complex) matrix, and R is upper triangular, with the
eigenvalues this time along the diagonal. The matrix A must be square.
Example 8.31
The Schur decomposition of the matrix we diagonalized in a previous example is as follows:
−5 4 9
A = −22 14 18 = Q · R · QT =
16 −8 −6
| {z }
A
−0.4082 0.1826 0.8944 −6 −20.1246 31.0376 −0.4082 −0.8165 0.4082
−0.8165 0.3651 −0.4472 0 3 5.7155 0.1826 0.3651 0.9129 .
0.4082 0.9129 0 0 0 6 0.8944 −0.4472 0
| {z }| {z }| {z }
Q R QT
If A is symmetric, then the upper triangular matrix R reduces to the diagonal matrix
with eigenvalues on the diagonal, Λ; the Schur decomposition is in this case simply A =
Q · Λ · QT .
where Qn×n and QH m×m are orthogonal (unitary, if complex) matrices, and B has positive
numbers µi , (i = 1, 2, . . . , r) in the first r positions on the main diagonal, and zero everywhere
else. It turns out that r is the rank of An×m . The columns of Qn×n are the eigenvectors of
An×m · AH H
n×m . The columns of Qm×m are the eigenvectors of An×m · An×m . The values µi ,
(i = 1, 2, . . . , r) ∈ R1 are called the singular values of A. They are analogous to eigenvalues
and are in fact are the positive square root of the eigenvalues of An×m ·AH H
n×m or An×m ·An×m .
Note that since the matrix from which the eigenvalues are drawn is Hermitian, that the
eigenvalues, and thus the singular values, are guaranteed real. Note also that if A itself is
4
Issai Schur, 1875-1941, Belrussian-born German-based mathematician.
square and Hermitian, that the absolute value of the eigenvalues of A will equal its singular
values. If A is square and non-Hermitian, there is no simple relation between its eigenvalues
and singular values. The factorization Qn×n · Bn×m · QH m×m is called the singular value
decomposition.
As discussed by Strang, the column vectors of Qn×n and Qm×m are even more than
orthonormal. They also must be chosen in such a way that An×m · Qm×m is a scalar multiple
of Qn×n . This comes directly from postmultiplying the general form of the singular value
decomposition by Qm×m : An×m · Qm×m = Qn×n · Bn×m . So in fact a more robust way
of computing the singular value decomposition is to first compute one of the orthogonal
matrices, and then compute the other orthogonal matrix with which the first one is consistent.
Example 8.32
Find the singular value decomposition of
1 −3 2
A2×3 = .
2 0 3
The matrix is real so we do not need to consider the conjugate transpose; we will retain the notation
for generality though here the ordinary transpose would suffice. First consider A · AH :
1 2
1 −3 2 14 8
A · AH = −3 0 = .
2 0 3 8 13
| {z } 2 3
A | {z }
AH
The diagonal eigenvalue matrix and corresponding orthogonal matrix composed of the normalized
eigenvectors in the columns are
21.5156 0 0.728827 −0.684698
Λ2×2 = , Q2×2 = .
0 5.48439 0.684698 0.728827
Next we consider AH · A:
1 2 5 −3 8
1 −3 2
AH · A = −3 0 = −3 9 −6 .
2 0 3
2 3 | {z } 8 −6 13
| {z } A
AH
The diagonal eigenvalue matrix and corresponding orthogonal matrix composed of the normalized
eigenvectors in the columns are
21.52 0 0 0.4524 0.3301 −0.8285
Λ3×3 = 0 5.484 0 , Q3×3 = −0.4714 0.8771 0.09206 .
0 0 0 0.7571 0.3489 0.5523
We take √
21.52 √ 0 0 4.639 0 0
B2×3 = = ,
0 5.484 0 0 2.342 0
It is also easily shown that the singular values of a square Hermitian matrix are identical
to the eigenvalues of that matrix. The singular values of a square non-Hermitian matrix are
not, in general, the eigenvalues of that matrix.
A = Q · H · QT ,
where Q is an orthogonal (or unitary) matrix and H has zeros below the first subdiagonal.
When A is Hermitian, Q is tridiagonal, which is very easy to invert numerically. Also H
has the same eigenvalues as A
Example 8.33
The Hessenberg form of our example square matrix A is
−5 4 9
A = −22 14 18 = Q · H · QT =
16 −8 −6
1 0 0 −5 2.0586 9.6313 1 0 0
0 −0.8087 0.5882 27.2029 2.3243 −24.0451 0 −0.8087 0.5882 .
0 0.5882 0.8087 0 1.9459 5.6757 0 0.5882 0.8087
| {z }| {z }| {z }
Q H QT
0 = (A · z)T · e,
= (A · z)T · (bp − b),
= zT · AT · (A · xp − b),
= zT · (AT · A · xp − AT · b).
AT · A · xp − AT · b = 0. (8.99)
from which
AT · A · xp = AT · b,
xp = (AT · A)−1 · AT · b,
A · xp = A · (AT · A)−1 · AT · b,
bp = A · (AT · A)−1 · AT ·b.
| {z }
≡R
R = A · (AT · A)−1 · AT .
The projection matrix for an operator A, when operating on an arbitrary vector b yields
the projection of b onto the column space of A. Note that many vectors b could have the
same projection onto the column space of A.
• substitutes each data point into the assumed form so as to form an overconstrained
system of linear equations,
• uses the technique associated with projection matrices to solve for the coefficients which
best represent the given data.
Example 8.34
Find the best straight line to approximate the measured data relating x to t.
t x
0 5
1 7
2 10
3 12
6 15
A straight line fit will have the form
x = a0 + a1 t,
where a0 and a1 are the terms to be determined. Substituting each data point to the assumed form,
we get five equations in two unknowns:
5 = a0 + 0a1 ,
7 = a0 + 1a1 ,
10 = a0 + 2a1 ,
12 = a0 + 3a1 ,
15 = a0 + 6a1 .
Rearranging, we get
1 0 5
1 1 7
a0
1 2 = 10 .
a1
1 3 12
1 6 15
This is of the form A · a = b. We then find that
−1
a = AT · A · AT · b.
Substituting, we find that
−1
1 0 5
1 1 7
a0 1 1 1 1 1 1 1 1 1 1 5.7925
= 1 2 10 = .
a1 0 1 2 3 6 0 1 2 3 6 1.6698
| {z }
| {z } 1 3 | {z } 12
a AT 1 6 AT 15
| {z } | {z }
A b
20
18
16
14
Data Points
12
10
x
6 x = 5.7925 + 1.6698 t
4
0
0 1 2 3 4 5 6 7
Figure 8.5: Plot of x − t data and best least squares straight line fit.
Example 8.35
Find the best straight line fit for the data in the previous example. Now however, assume that we
have five times the confidence in the accuracy of the final two data points, relative to the other points.
Define a square weighting matrix W:
1 0 0 0 0
0 1 0 0 0
W = 0 0 1 0 0.
0 0 0 5 0
0 0 0 0 5
Now we perform the following operations:
A·a = b,
W·A·a = W · b,
T T
(W · A) · W · A · a = (W · A) · W · b,
−1
a = (W · A)T · W · A (W · A)T · W · b.
20
18
weighted data points
16
14
12
10
x
8
x = 8.0008 + 1.1972 t
6
0
0 1 2 3 4 5 6 7
Figure 8.6: Plot of x − t data and best weighted least squares straight line fit.
x = 8.0008 + 1.1972 t.
A plot of the raw data and the best fit straight line is shown in Figure 8.6
When the measurements are independent and equally reliable, W is the identity matrix.
If the measurements are independent but not equally reliable, W is at most diagonal. If the
measurements are not independent, then non-zero terms can appear off the diagonal in W.
It is often advantageous, for instance in problems in which one wants to control a process in
real time, to give priority to recent data estimates over old data estimates and to continually
employ a least squares technique to estimate future system behavior. The previous example
does just that. A famous fast algorithm for such problems is known as a Kalman Filter.
d At 1
e = A + A2 t + A3 t2 + · · · ,
dt 2!
1 22 1 33
= A · I + At + A t + A t + · · · ,
2! 3!
| {z }
=eAt
At
= A·e .
eaI = ea I, (8.101)
(eA )−1 = e−A , (8.102)
eA(t+s) = eAt eAs . (8.103)
Example 8.36
Find eAt if
a 1 0
A= 0 a 1 .
0 0 a
We have
A = aI + B,
where
0 1 0
B= 0 0 1 .
0 0 0
Thus
0 0 1
B2 = 0 0 0 ,
0 0 0
0 0 0
B3 = 0 0 0 ,
0 0 0
..
.
0 0 0
Bn = 0 0 0 , for n ≥ 4.
0 0 0
Furthermore
I · B = B · I = B.
Thus
eAt = e(aI+B)t ,
= eatI · eBt ,
=0
z }| {
I + atI + 1 a2 t2 I2 + 1 a3 t3 I3 + · · · · I + Bt + 1 B2 t2 + 1 B3 t3 + · · · ,
=
| 2! {z 3! }
| 2! {z 3! }
=eatI =eat I =eBt
2
at 2t
= e I · I + Bt + B ,
2
2
1 t t2 ,
= eat 0 1 t .
0 0 1
it is not obvious whether or not there exist (x1 , x2 , x3 ) which will give positive or negative
values of f . However, it is easily verified that f can be rewritten as
So in this case f ≥ 0 for all (x1 , x2 , x3 ). How to demonstrate positivity (or non-positivity)
of such expressions is the topic of this section. A quadratic form is an expression
n X
X n
f (x1 , · · · , xn ) = aij xi xj , (8.104)
j=1 i=1
Example 8.37
Change
f (x1 , x2 ) = 2x21 + 2x1 x2 + 2x22 ,
to standard form.
For n = 2, equation (8.104) becomes
a11 = 2,
a12 = 1,
a21 = 1,
a22 = 2.
So we get
2 1
A= .
1 2
The eigenvalue of A are λ = 1, λ = 3. The orthogonal matrix corresponding to A is
1 1 1 −1 T 1 1 −1
Q= √ , Q =Q = √ .
2 −1 1 2 1 1
The transformation x = Q · y is
1
x1 = √ (y1 + y2 ),
2
1
x2 = √ (−y1 + y2 ).
2
Example 8.38
Change
f (x1 , x2 , x3 ) = 18x21 − 16x1 x2 + 5x22 + 12x1 x3 − 4x2 x3 + 6x23 ,
to standard form.
It is clear that f (x1 , x2 , x3 ) is positive definite. Moreover, by carrying out the multiplications, it is easily
seen that the original form is recovered. Further manipulation would also show that f (x1 , x2 , x3 ) =
2(x1 − x2 + x3 )2 + 3(2x1 − x2 )2 + 4(x1 + x3 )2 , so we see the particular quadratic form is not unique.
An×m · A+
m×n · An×m = An×m . (8.114)
5
after Eliakim Hastings Moore, 1862-1932, American mathematician, and Sir Roger Penrose, 1931-, En-
glish mathematician. It is also credited to Arne Bjerhammar, 1917-, Swedish geodesist.
A+ + H
m×n = Qm×m · Bm×n · Qn×n . (8.115)
An×m · A+
m×n = In×n . (8.116)
Let’s check this with our definitions for the case when n ≤ m, n = r.
An×m · A+ m×n = Q n×n · B n×m · QH
m×m · Qm×m · B +
m×n · QH
n×n , (8.117)
= Qn×n · Bn×m · Q−1
m×m · Qm×m · B+
m×n · QH
n×n , (8.118)
= Qn×n · Bn×m · B+ H
m×n · Qn×n , (8.119)
= Qn×n · In×n · QH
n×n , (8.120)
H
= Qn×n · Qn×n , (8.121)
= Qn×n · Q−1
n×n , (8.122)
= In×n . (8.123)
We note for this special case that precisely because of the way we defined B+ that Bn×m ·
B+ +
m×n = In×n . When n > m, Bn×m · Bm×n yields a matrix with r ones on the diagonal and
zeros elsewhere.
Example 8.39
Find the Moore-Penrose inverse, A+
3×2 , of A2×3 in the previous example:
1 −3 2
A2×3 = .
2 0 3
A+
3×2 = Q3×3 · B+ H
3×2 · Q2×2 ,
1
0.452350 0.330059 −0.828517 4.6385 0
0.728827 0.684698
A+
3×2 = −0.471378 0.877114 0.0920575 0 1
2.3419
,
−0.684698 0.728827
0.757088 0.348902 0.552345 0 0
−0.0254237 0.169492
A+
3×2 = −0.330508 0.20339 .
0.0169492 0.220339
Note that
−0.0254237 0.169492
1 −3 2 −0.330508 0.20339 = 1 0 .
A2×3 A+
3×2 = 2 0 3 0 1
0.0169492 0.220339
Example 8.40
Use the Moore-Penrose inverse to solve the problem A · x = b studied in an earlier example:
1 2 x1 2
= .
3 6 x2 0
We will need B+ , which is easily calculated by taking the inverse of each diagonal term of B:
1
√ 0
B+ = 5 2 .
0 0
We see that the Moore-Penrose operator acting on b has yielded an x vector which is in the row space
of A. As there is no right null space component, it is the minimum length vector that minimizes the
error ||A · x − b||2 . It is fully consistent with the solution we found using Gaussian elimination in an
earlier example.
Problems
1. Find the x with smallest ||x||2 which minimizes ||A · x − b||2 for
1 0 3 1
A = 2 −1 3 , b = 0.
3 −1 5 1
3. Find x with the smallest ||x||2 which minimizes ||A · x − b||2 for
1 0 1 4 2
A = 1 0 2 −1 , b = 1 .
2 1 3 −2 −3
4. Find eA if
1 1 1
A= 0 3 2 .
0 0 5
5. Diagonalize or reduce to Jordan canonical form
5 2 −1
A= 0 5 1 .
0 0 5
7. Decompose A into Jordan form S · J · S−1 , P−1 · L · D · U, Q · R, Schur form, and Hessenberg form
0 1 0 1
1 0 1 0
A= 0 1 0 1 .
1 0 1 0
8. Find the matrix S that will convert the following to the Jordan canonical form
6 −1 −3 1
−1 6 1 −3
(a) ,
−3 1 6 −1
1 −3 −1 6
8 −2 −2 0
0 6 2 −4
(b) −2 0
,
8 −2
2 −4 0 6
and show the Jordan canonical form.
1 1 2 0
0 1 3 0
9. Show that the eigenvectors and generalized eigenvectors of
span the space.
0 0 2 2
0 0 0 1
10. Find the projection matrix onto the space spanned by (1, 1, 1) and (1, 2, 3).
11. Reduce 4x2 + 4y 2 + 2z 2 − 4xy + 4yz + 4zx to standard quadratic form.
12. Find the inverse of
1/4 1/2 3/4
3/4 1/2 1/4 .
1/4 1/2 1/2
0 0 i
13. Find exp 0 1 0 .
1 0 0
1 3
14. Find the nth power of .
3 1
15. If
5 4
A= ,
1 2
find a matrix S such that S−1 · A · S is a diagonal matrix. Show by multiplication that it is indeed
diagonal.
6 2 8 6
16. Determine if A = and B = are similar.
−2 1 −3 −1
17. Find the eigenvalues, eigenvectors, and the matrix S such that S−1 · A · S is diagonal or of Jordan
form, where A is
5 0 0
(a) 1 0 1 ,
0 0 −2
−2 0 2
(b) 2 1 0 ,
0 0 −2i
3 0 −1
(c) −1 2 2i .
1 0 1+i
Dynamical systems
In this chapter we consider the evolution of systems, often called dynamic systems. Generally,
we will be concerned with systems which can be described by sets of ordinary differential
equations, both linear and non-linear. Some other classes of systems will also be studied.
• If not linearizable, attempt to ascertain the stability of the non-linear system near its
equilibria.
353
354 CHAPTER 9. DYNAMICAL SYSTEMS
Example 9.1
For x ∈ R2 , t ∈ R1 , f : R2 → R2 , consider
dx1
= x2 − x21 = f1 (x1 , x2 ),
dt
dx2
= x2 − x1 = f2 (x1 , x2 ).
dt
The curves defined in the x1 , x2 plane by f1 = 0 and f2 = 0 are very useful in determining both the fixed
points (found at the intersection) and in the behavior of the system of differential equations. In fact
one can sketch trajectories of paths in this phase space by inspection in many cases. The loci of points
where f1 = 0 and f2 = 0 are plotted in Figure 9.1. The zeroes are found at (x1 , x2 )T = (0, 0)T , (1, 1)T .
Linearize about both points to find the local behavior of the solution near these points. Near (0,0), the
linearization is
dx1
= x2 ,
dt
dx2
= x2 − x1 ,
dt
or
d x1 0 1 x1
= .
dt x2 −1 1 x2
This is of the form
dx
= A · x.
dt
And with
P · z ≡ x,
where P is a constant matrix, we get
d dz
(P · z) = P · = A · P · z,
dt dt
dz
= P−1 · A · P · z.
dt
At this point we assume that A has distinct eigenvalues and linearly independent eigenvectors; other
cases are easily handled. If we choose P such that its columns contain the eigenvectors of A, we will
get a diagonal matrix, which will lead to a set of uncoupled differential equations; each of these can be
solved individually. So for our A, standard linear algebra gives
√ √ √ !
√i 1
+ 63 i
1 3 1 3
+ i − i −1 3 2
P= 2 2 2 2 , P = √ .
1 1 − √i3 12 − 63 i
x2 f1= 0 f2= 0
1.5
0.5
0 x1
−0.5
√ !
dz2 1 3
= + i z2 ,
dt 2 2
| {z }
=λ2
Since there is a positive real coefficient in the exponential terms, both x1 and x2 grow exponentially.
The imaginary component indicates that this is an oscillatory growth. Hence, there is no tendency for
a solution which is initially close to (0, 0), to remain there. So the fixed point is unstable.
Consider the next fixed point near (1, 1). First define a new set of local variables:
x̃1 = x1 − 1,
x̃2 = x2 − 1.
Then
dx1 dx̃1
= = (x̃2 + 1) − (x̃1 + 1)2 ,
dt dt
dx2 dx̃2
= = (x̃2 + 1) − (x̃1 + 1).
dt dt
Expanding, we get
dx̃1
= (x̃2 + 1) − x̃21 − 2x̃1 − 1,
dt
dx̃2
= (x̃2 + 1) − (x̃1 + 1).
dt
Linearizing about (x̃1 , x̃2 ) = (0, 0), we find
dx̃1
= x̃2 − 2x̃1 ,
dt
dx̃2
= x̃2 − x̃1 ,
dt
or
d x̃1 −2 1 x̃1
= .
dt x̃2 −1 1 x̃2
Going through an essentially identical exercise gives the eigenvalues to be
√
1 5
λ1 = − + > 0,
2 2
√
1 5
λ2 = − − < 0,
2 2
which in itself shows the solution to be essentially unstable since there is a positive eigenvalue. After
the usual linear algebra and back transformations, one obtains the local solution:
√ ! √ ! ! √ ! √ ! !
3− 5 1 5 3+ 5 1 5
x1 = 1 + c1 exp − + t + c2 exp − − t ,
2 2 2 2 2 2
√ ! ! √ ! !
1 5 1 5
x2 = 1 + c1 exp − + t + c2 exp − − t .
2 2 2 2
Note that while this solution is generally unstable, if one has the special case in which c1 = 0, that the
fixed point in fact is stable. Such is characteristic of a saddle node.
Example 9.2
For x ∈ R2 , t ∈ R1 , f : R2 × R1 → R2 , analyze
dx1 dx2
t + x2 x1 = x1 + t = f1 (x1 , x2 , t),
dt dt
dx1 dx2
x1 + x22 = x1 t = f2 (x1 , x2 , t),
dt dt
x1 (0) = x10 , x2 (0) = x20 .
Let
dt
= 1, t(0) = 0,
ds
and further y1 = x1 , y2 = x2 , y3 = t. Then with s ∈ R1 , y ∈ R3 , g : R3 → R3 ,
dy1 dy2
y3 + y2 y1 = y1 + y3 = g1 (y1 , y2 , y3 ),
ds ds
dy1 2 dy2
y1 + y2 = y1 y3 = g2 (y1 , y2 , y3 ),
ds ds
dy3
= 1 = g3 (y1 , y2 , y3 ),
ds
y1 (0) = y10 , y2 (0) = y20 , y3 (0) = 0.
Inverting the coefficient matrix, we obtain the following equation which is in autonomous form:
y y −y2 y +y y
1 2 1 3 2 3
y y2 y3 −y12 h1 (y1 , y2 , y3 )
d 1 y1 (y −y1 −y3 )
2
y2 = y y3 y −y2 = h2 (y1 , y2 , y3 ) .
ds 2( 2 3 )
y3 1 h3 (y1 , y2 , y3 )
1
There are potential singularities at y2 = 0 and y2 y3 = y12 . These can be addressed by defining a new
independent variable u ∈ R1 via the equation
ds
= y2 y2 y3 − y12 .
du
The system of equations then transforms to
y y y y − y 2 y + y2 y3 p1 (y1 , y2 , y3 )
d 1 2 1 22 1 3
y2 = y1 y3 − y1 − y3 = p2 (y1 , y2 , y3 ) .
du 2
y3 y2 y2 y3 − y1 p3 (y1 , y2 , y3 )
This equation actually has an infinite number of fixed points, all of which lie on a line in the three
dimensional phase volume. The line is given parametrically by (y1 , y2 , y3 )T = (0, 0, v)T , v ∈ R1 Here
v is just a parameter used in describing the line of fixed points. However, it turns out in this case
that the Taylor series expansions yield no linear contribution any of the fixed points, so we don’t get
to use the standard linear analysis technique! The problem has an essential non-linear essence, even
near fixed points. More potent methods would need to be employed, but the example demonstrates the
principle. Figure 9.2 gives a numerically obtained solution for y1 (u), y2 (u), y3 (u) along with a trajectory
in y1 , y2 , y3 space when y1 (0) = 1, y2 (0) = −1, y3(0) = 0. This corresponds to x1 (t = 0) = 1, x2 (t =
0) = −1.
We note that while the solutions are monotonic in the variable u, that they are not monotonic
in t, after the transformation back to x1 (t), x2 (t) is effected. Also, while it appears there are points
(u = 0.38, u = 0.84, u = 1.07) where the derivatives dy 1 dy2 dy3
du , du , du become unbounded, closer inspection
reveals that they are simply points of steep, but bounded, derivatives. However at points where the
slop dy dt dx1 dx2
du = du changes sign, the derivatives dt , dt formally are infinite, as is reflected in the cyclic
3
2
y 1 y = x
2 1 1
0
-1 50
-2
10 40
8 30
6 20
y
3 4
10
2
u
0 0.2 0.4 0.6 0.8 1 1.2
0
y =x
20 2 2
x 2
1 y 40
1
1
50
u
40 0.2 0.4 0.6 0.8 1 1.2
-1
30
-2
20
y =t
3
10 t 10
2 4 6 8 10
x2
8
2
6
1 4
2
t
2 4 6 8 10
u
0.2 0.4 0.6 0.8 1 1.2
-1
-2
Figure 9.2: Solutions for one set of initial conditions, y1 (0) = 1, y2 (0) = −1, y3 (0) = 0,
for second paradigm example: trajectory in phase volume (y1 , y2, y3 ); also y1 (u), y2 (u), y3(u)
and x1 (t), x2 (t). Here y1 = x1 , y2 = x2 , y3 = t.
g1 (y1 , . . . , yn+1) f1 (x1 , . . . , xn , t)
.. ..
g(y) = . .
g (y , . . . , y ) f (x , . . . , x , t) .
= (9.6)
n 1 n+1 n 1 n
gn+1 (y1 , . . . , yn+1) 1
The original equation then is of the form
dy
B(y) · = g(y). (9.7)
ds
By forming B−1 , it can be written as
dy
= B−1 (y) · g(y), (9.8)
ds
or by taking
B−1 (y) · g(y) ≡ h(y), (9.9)
we get the form, commonly called autonomous form, with s ∈ R1 , y ∈ Rn+1 , h : Rn+1 →
Rn+1 :
dy
= h(y). (9.10)
ds
Sometimes h has singularities. If the source of the singularity can be identified, a
singularity-free autonomous set of equations can often be written. For example, suppose
h can be rewritten as
p(y)
h(y) = (9.11)
q(y)
where p and q have no singularities. Then we can remove the singularity by introducing the
new independent variable u ∈ R1 such that
ds
= q(y). (9.12)
du
Using the chain rule, the system then becomes
dy p(y)
= , (9.13)
ds q(y)
ds dy p(y)
= q(y) , (9.14)
du ds q(y)
dy
= p(y), (9.15)
du
(9.16)
• Find all the zeroes of h. This is an algebra problem, which can be topologically difficult
for non-linear problems.
• If h has any singularities, redefine variables in the manner demonstrated to remove the
singularity
• If the system is linear, an eigenvalue analysis is sufficient to reveal stability; for non-
linear systems, the situation is not always straightforward.
xk+1
i = fi (xk1 , xk2 , · · · , xkn ), i = 1, · · · , n. (9.17)
Given an initial point x0i , (i = 1, . . . , n) in Rn , a series of images x1i , x2i , x3i , . . . can be found
as k = 0, 1, 2, . . .. The map is dissipative or conservative according to whether the diameter
of a set is larger than that of its image or the same, respectively, i.e if the determinant of
the Jacobian matrix ∂fi /∂xj is < or = 1.
The point xi = xi is a fixed point of the map if it maps to itself, i.e. if
xi = fi (x1 , x2 , · · · , xn ), i = 1, · · · , n. (9.18)
The fixed point xi = 0 is linearly unstable if a small perturbation from it leads the images
farther and farther away. Otherwise it is stable. A special case of this is asymptotic stability
wherein the image returns arbitrarily close Pto the fixed point.
A linear map can be written as xi = nj=1 Aij xkj , (i = 1, 2, . . .) or xk+1 = A · xk . The
k+1
origin x = 0 is a fixed point of this map. If ||A|| > 1, then ||xk+1 || > ||xk || and the map is
unstable. Otherwise it is stable.
Example 9.3
Examine the linear stability of the fixed points of the logistics map, popularized by May.1
xk+1 = rxk (1 − xk ),
We take r ∈ [0, 4] so that xk ∈ [0, 1] maps onto xk+1 ∈ [0, 1]. That is the mapping is onto itself.
The fixed points are solutions of
x = rx(1 − x),
1
Robert McCredie May, 1936-, Australian-Anglo ecologist.
which are
1
x = 0, x=1− .
r
Consider the mapping itself. For an initial seed x0 , we generate a series of xk . For example if we take
r = 0.4 and xo = 0.3, we get
x0 = 0.3,
x1 = 0.4(0.3)(1 − 0.3) = 0.084,
x2 = 0.4(0.084)(1 − 0.084) = 0.0307776,
x3 = 0.4(0.0307776)(1 − 0.0307776) = 0.0119321,
x4 = 0.4(0.0119321)(1 − 0.0119321) = 0.0047159,
x5 = 0.4(0.0047159)(1 − 0.0047159) = 0.00187747,
..
.
x∞ = 0.
4
For this value of r, the solution approaches the fixed point of 0. Consider r = 3 and x0 = 0.3
x0 = 0.3,
x1 = (4/3)(0.3)(1 − 0.3) = 0.28,
x2 = (4/3)(0.28)(1 − 0.28) = 0.2688,
x3 = (4/3)(0.2688)(1 − 0.2688) = 0.262062,
x4 = (4/3)(0.262062)(1 − 0.262062) = 0.257847,
x5 = (4/3)(0.257847)(1 − 0.257847) = 0.255149,
..
.
1
x∞ = 0.250 = 1 − .
r
In this case the solution was attracted to the alternate fixed point.
To analyze the stability of each fixed point, we give it a small perturbation x′ . Thus x + x′ is
mapped to x + x′′ , where
Simplifying, we get
x′′ = rx′ (1 − 2x).
A fixed point is stable if |x′′ /x′ | ≤ 1. This indicates that the perturbation is decaying. Now consider
each fixed point in turn.
x = 0:
x
1.0
0.8
0.6
0.4
0.2
r
1 2 3 4
Figure 9.3: Plot of xk as k → ∞ as a function of r for the logistics map, xk+1 = rxk (1 − xk )
for r ∈ [0, 4].
dn x dn−1 x dx
n
+ an (x, t) n−1
+ · · · + a2 (x, t) + a1 (x, t)x = f (t), (9.19)
dt dt dt
dy1
= y2 , (9.20)
ds
dy2
= y3 , (9.21)
ds
..
. (9.22)
dyn−1
= yn , (9.23)
ds
dyn
= −an (y1 , yn+1)yn − an−1 (y1, yn+1 )yn−1 − · · · − a1 (y1 , yn+1)y1 + f (yn+1 ),(9.24)
ds
dyn+1
= 1. (9.25)
ds
CC BY-NC-ND. 28 March 2011, M. Sen, J. M. Powers.
9.4. HIGH ORDER SCALAR DIFFERENTIAL EQUATIONS 365
Example 9.4
For x ∈ R1 , t ∈ R1 , consider the forced Duffing equation:
d2 x dx
+ x + x3 = sin(2t), x(0) = 0, = 0.
dt2 dt t=0
Here a2 (x, t) = 0, a1 (x, t) = 1 + x2 , f (t) = sin(2t). Now this non-linear differential equation with
homogeneous boundary conditions and forcing has no analytic solution. It can be solved numerically;
most solution techniques require a recasting as a system of first order equations. To recast this as an
autonomous set of equations, with y ∈ R3 , s ∈ R1 , consider
dx
x = y1 , = y2 , t = s = y3 .
dt
d d
Then dt = ds , and the equations transform to
y y2 h1 (y1 , y2 , y3 ) y1 (0) 0
d 1
y2 = −y1 − y13 + sin(2y3 ) = h2 (y1 , y2 , y3 ) , y2 (0) = 0 .
ds
y3 1 h3 (y1 , y2 , y3 ) y3 (0) 0
Note that this system has no equilibrium point as there exists no y for which h = 0. Once the numerical
solution is obtained, one transforms back to x, t space. Figure 9.4 give the trajectory in the y1 , y2 , y3
phase space, and a plot of the corresponding solution x(t) for t ∈ [0, 50].
-1 -0.5 y1 = x
y = dx/dt 0
2 1 0.5
0 1
-1
40
0.5
t
10 20 30 40 50
y3 = t
-0.5
20
-1
Figure 9.4: Phase space trajectory and solution x(t) for forced Duffing equation.
9.5.1.1 n eigenvectors
We will assume that there is a full set of eigenvectors even though not all the eigenvalues
are distinct. If e1 , e2 , . . . , en are the eigenvectors corresponding to eigenvalues λ1 , λ2 , . . . , λn ,
then n
X
x= ci ei eλi t , (9.34)
i=1
is the general solution, where c1 , c2 , . . . , cn are arbitrary constants.
Example 9.5
For x ∈ R3 , t ∈ R1 , A ∈ R3 × R3 , solve dx
dt = A · x where
1 −1 4
A = 3 2 −1 .
2 1 −1
The eigenvalues and eigenvectors are
−1
λ1 = 1, e1 = 4 ,
1
1
λ2 = 3, e2 = 2 ,
1
−1
λ3 = −2, e3 = 1 .
1
Thus the solution is
−1 1 −1
x = c1 4 et + c2 2 e3t + c3 1 e−2t .
1 1 1
or expanding,
x1 (t) = −c1 et + c2 e3t − c3 e−2t ,
x2 (t) = 4c1 et + 2c2 e3t + c3 e−2t ,
x3 (t) = c1 et + c2 e3t + c3 e−2t .
Example 9.6
For x ∈ R3 , t ∈ R1 , A ∈ R3 × R3 , solve dx
dt = A · x where
2 −1 −1
A = 2 1 −1 .
0 −1 1
One solution of dx
dt
= A · x is x = eAt · c, where c is a constant vector. If c1 , c2 ,· · ·, cn
are linearly independent vectors, then xi = eAt · ci , i = 1, · · · , n, are linearly independent
solutions. We would like to choose ci , i = 1, 2, · · · , n, such that each eAt · ci is a series with
a finite number of terms. This can be done in the following manner. Since
If we apply the matrix operator S−1 , which is a constant, to both sides, we get
!
d
S−1 · x = J · |S−1
{z· x} . (9.41)
dt | {z }
≡z ≡z
Example 9.7
For x ∈ R3 , t ∈ R1 , A ∈ R3 × R3 , find the general solution of
dx
= A · x,
dt
where
4 1 3
A= 0 4 1 .
0 0 4
A has an eigenvalue λ = 4 with multiplicity three. The eigenvector is
1
e = 0 ,
0
Alternative method
Alternatively, we can simply use the Jordan decomposition to form the solution. When we form
the matrix S from the eigenvectors and generalized eigenvectors, we have
1 0 0
S = ( e g1 g2 ) = 0 1 −3 .
0 0 1
We then get
1 0 0
S−1 = 0 1 3 ,
0 0 1
4 1 0
J = S−1 · A · S = 0 4 1 .
0 0 4
dz
Now with z = S−1 · x, we solve dt = J · z,
z 4 1 0 z1
d 1
z2 = 0 4 1 z2 .
dt
z3 0 0 4 z3
dz3
The final equation is totally uncoupled; solving dt = 4z3 , we get
z3 (t) = c3 e4t .
x = Ω · c, (9.44)
where
c1
c = ... .
(9.45)
cn
The term eAt = Ω(t) · Ω−1 (0) is a fundamental matrix.
Example 9.8
Find the fundamental matrix of the problem given above.
so that
t2 c1
1 t 2
x = Ω · c = e4t 0 1 −3 + t c2 .
0 0 1 c3
Example 9.9
For x ∈ R2 , t ∈ R1 , solve
dx1
= 2x1 + x2 + 1,
dt
dx2
= x1 + 2x2 + t.
dt
This can be written as
d x1 2 1 x1 1
= + .
dt x2 1 2 x2 t
We have
1 1 −1 1 1 −1 1 0
P= , P = , Λ= ,
−1 1 2 1 1 0 3
so that
d z1 1 0 z1 1 1−t
= + .
dt z2 0 3 z2 2 1+t
The solution is
t
z1 = aet + ,
2
2 t
z2 = be3t − − ,
9 6
Example 9.10
Solve the system
dx
= A · (x − xo ) + b, x(to ) = xo .
dt
Such a system arises naturally when one linearizes a non-linear system of the form dx/dt = f (x) about
a point x = xo . Here then, A is the Jacobian matrix A = ∂f /∂x|x=xo . Note that the system is in
equilibrium when
A · (x − xo ) = −b,
or
x = x0 − A−1 · b.
Further note that if b = 0 that the initial condition x = xo is also an equilibrium condition, and is the
unique solution to the differential equation.
First define a new dependent variable z:
z ≡ x − xo + A−1 · b.
So we have
x = z + xo − A−1 · b.
At t = to , we then get
z(to ) = A−1 · b.
Then using the definition of z, one can write the solution in terms of the original x as
x(t) = xo + P · eΛ(t−to ) · P−1 − I · A−1 · b.
Note that the time scales of evolution are entirely determined by Λ; in particular the time scales of
each mode, τi , are τi = 1/λi , where λi is an entry in Λ. The constant vector b plays a secondary role
in determining the time scales.
Lastly, one recalls from the definition of the matrix exponential that eA(t−to ) = P · eΛ(t−to ) · P−1 ,
so we get the final form or
x(t) = xo + eA(t−to ) − I · A−1 · b.
Example 9.11
dx
For x ∈ R3 , t ∈ R1 , A ∈ R3 × R3 , f : R1 → R3 , solve dt = A · x + f (t) with
4 1 3 3et
A= 0 4 1 , f = 0 .
0 0 4 0
The homogeneous part of this problem has been solved before. Let the particular solution be
xP = cet .
Therefore,
−1
x = xH + 0 et .
0
This follows the general procedure explained in Section 3.3.2, page 79.
9.6.1 Definitions
With x ∈ Rn , t ∈ R1 , f : Rn → Rn , consider a system of n nonlinear first-order ordinary
differential equations
dxi
= fi (x1 , x2 , · · · , xn ), i = 1, · · · , n. (9.50)
dt
where t is time, and fi is a vector field. The system is autonomous since fi is not a function
of t. The coordinates
Pn x1 , x2 , · · · , xn form a phase or state space. The divergence of the
vector field i=1 ∂fi /∂xi indicates the change of a given volume of initial conditions in
phase space. If the divergence is zero, the volume remains constant and the system is said to
be conservative. If the divergence is negative, the volume shrinks with time and the system
is dissipative. The volume in a dissipative system eventually goes to zero. This final state
to which some initial set of points in phase space goes is called an attractor. Attractors may
be points, closed curves, or tori, or fractal (strange). A given dynamical system may have
several attractors that co-exist. Each attractor has its own basin of attraction in Rn ; initial
conditions that lie on this basin tend to that particular attractor.
The steady state solutions xi = xi of equation (9.50) are called critical (or fixed, singular
or stationary) points. Thus, by definition
fi (x1 , x2 , · · · , xn ) = 0, i = 1, · · · , n, (9.51)
which is an algebraic or transcendental equation. The dynamics of the system is analyzed
by studying the stability of the critical point. For this we perturb the system so that
xi = xi + x′i , (9.52)
where the prime denotes a perturbation. If ||x′i || is bounded for t → ∞, the critical point
is said to be stable, otherwise it is unstable. As a special case, if ||x′i || → 0 as t → ∞, the
critical point is asymptotically stable.
into (9.50), and linearize by keeping only the terms that are linear in x′i and neglecting all
products of x′i . Thus equation (9.50) takes a linearized local form
n
dx′i X
= Aij x′j . (9.53)
dt j=1
Another way of obtaining the same result is to expand the vector field in a Taylor series
around xi = xi so that
n
X ∂fi
fi (xi ) = x′j + H.O.T., (9.54)
j=1
∂xj xi =xi
and then neglect the higher order terms (H.O.T.) Thus, in equation (9.53)
∂fi
Aij = (9.55)
∂xj xi =xi
is the Jacobian of fi evaluated at the critical point. In matrix form the linearized equation
for the perturbation x′ is
dx′
= A · x′ . (9.56)
dt
The real parts of the eigenvalues of A determine the linear stability of the critical point
x = 0, and the behavior of the solution near it:
1. If all eigenvalues have real parts < 0, the critical point is asymptotically stable.
2. If at least one eigenvalue has a real part > 0, the critical point is unstable.
3. If all eigenvalues have real parts ≤ 0, and some have zero real parts, then the critical
point is stable if A has k linearly independent eigenvectors for each eigenvalue of
multiplicity k. Otherwise it is unstable.
The following are some terms used in classifying critical points according to the real and
imaginary parts of the eigenvalues of A.
Classification Eigenvalues
Hyperbolic Non-zero real part
Saddle Some real parts negative, others positive
Stable node or sink All real parts negative
ordinary sink All real parts negative, imaginary parts zero
spiral sink All real parts negative, imaginary parts non-zero
Unstable node or source All real parts positive
ordinary source All real parts positive, imaginary parts zero
spiral source All real parts positive, imaginary parts non-zero
Center All purely imaginary and non-zero
Figures 9.5 and 9.6 show examples of phase planes for simple systems which describe
an ordinary source node, a spiral sink node, an ordinary center node, and a saddle node.
Figure 9.7 gives a phase plane, vector field, and trajectories for a complex system with many
nodes present. Here the nodes are spiral and saddle nodes.
1.5
0.5
0 y’=0
y
−0.5
−1
−1.5
−2
−3 −2 −1 0 1 2 3
x
1.5
0.5
0
y
−0.5
−1
−1.5
−2
−3 −2 −1 0 1 2 3
x
Figure 9.5: Phase plane for system with ordinary source node and spiral sink node.
• V > 0 for xi 6= 0,
• V = 0 for xi = 0,
dV
• dt
< 0 for xi 6= 0, and
1.5
0.5
0
y
−0.5
−1
−1.5
−2
−3 −2 −1 0 1 2 3
x
x’=y−x
y’=x
Saddle Node
y’=0 x’ = 0
2
1.5
0.5
0
y
−0.5
−1
−1.5
−2
−3 −2 −1 0 1 2 3
x
Figure 9.6: Phase plane for systems with center node and saddle node
dV
• dt
= 0 for xi = 0,
then the equilibrium point of the differential equations, xi = 0, is globally stable to all per-
turbations, large or small. The function V (x1 , x2 , · · · , xn ) is called a Lyapunov3 function.
Although one cannot always find a Lyapunov function for a given system of differential
equations, we can pose a method to seek a Lyapunov function given a set of autonomous
ordinary differential equations. While the method lacks robustness, it is always straight-
forward to guess a functional form for a Lyapunov function and test whether or not the
proposed function satisfies the criteria:
2
x’ = 0
1.5
x’ = 0
1
0.5
0
y
x’ = 0
−0.5
−1
−1.5
x’ = 0
−2
−3 −2 −1 0 1 2 3
y’ = 0 x y’ = 0
Example 9.12
Show that x = 0 is globally stable, if
d2 x dx
m 2
+β + k1 x + k2 x3 = 0, where m, β, k1 , k2 > 0.
dt dt
This system models the motion of a mass-spring-damper system when the spring is non-linear. Breaking
the original second order differential equation into two first order equations, we get
dx
= y,
dt
dy β k1 k2
= − y − x − x3 .
dt m m m
Here x represents the position, and y represents the velocity. Let us guess that the Lyapunov function
has the form
V (x, y) = ax2 + by 2 + cx4 , where a, b, c > 0.
Note that V (x, y) ≥ 0 and that V (0, 0) = 0. Then
dV ∂V dx ∂V dy
= + ,
dt ∂x dt ∂y dt
dx dx dy
= 2ax + 4cx3 + 2by ,
dt dt dt
β k1 k2
= (2ax + 4cx )y + 2by − y − x − x3 ,
3
m m m
bk1 bk2 2b
= 2 a− xy + 2 2c − x3 y − βy 2 .
m m m
k2
If we choose b = m 1
2 , a = 2 k1 , c = 4 , then the coefficients on xy and x3 y in the expression for dV
dt are
identically zero, and we get
dV
= −βy 2 ,
dt
which for β > 0 is negative for all y 6= 0 and zero for y = 0. Further, with these choices of a, b, c, the
Lyapunov function itself is
1 1 1
V = k1 x2 + k2 x4 + my 2 ≥ 0.
2 4 2
Checking, we see
dV dx dx dy
= k1 x + k2 x3 + my ,
dt dt dt dt
β k1 k2
= k1 xy + k2 x y + my − y − x − x3 ,
3
m m m
= k1 xy + k2 x3 y − βy 2 − k1 xy − k2 x3 y,
= −βy 2 ≤ 0.
Thus V is a Lyapunov function, and x = y = 0 is globally stable. Actually, in this case, V = (kinetic
energy + potential energy), where kinetic energy = 12 my 2 , and potential energy = 21 k1 x2 + 14 k2 x4 .
Note that V (x, y) is just an algebraic function of the system’s state variables. When we take the time
derivative of V , we are forced to invoke our original system, which defines the differential equations.
We note for this system that precisely since V is strictly positive or zero for all x, y, and moreover that
it is decaying for all time, that this necessarily implies that V → 0, hence x, y → 0.
dH ∂H dxi ∂H dyi
= + = 0, (9.61)
dt ∂xi dt ∂yi dt
dH ∂H ∂H
= fi (x1 , · · · , xn , y1 , · · · , yn ) + gi (x1 , · · · , xn , y1 , · · · , yn ) = 0. (9.62)
dt ∂xi ∂yi
This differential equation can at times be solved directly by the method of separation of
variables in which we assume a specific functional form for H(xi , yi).
Alternatively, we can also determine H by demanding that
∂H dxi ∂H dyi
= , =− . (9.63)
∂yi dt ∂xi dt
Substituting from the original differential equations, we are led to equations for H(xi , yi )
∂H ∂H
= fi (x1 , · · · , xn , y1, · · · , yn ), = −gi (x1 , · · · , xn , y1 , · · · , yn ). (9.64)
∂yi ∂xi
Example 9.13
Find the Hamiltonian for a linear mass spring system:
d2 x dx
m + kx = 0, x(0) = xo , = ẋ0 .
dt2 dt 0
dx
Taking dt = y to reduce this to a system of two first order equations, we have
dx
= y, x(0) = xo
dt
dy k
= − x, y(0) = yo
dt m
For this system n = 1.
dH
We seek H(x, y) such that dt = 0. That is
dH ∂H dx ∂H dy
= + = 0.
dt ∂x dt ∂y dt
As with all partial differential equations, one has to transform to a system of ordinary equations in
order to solve. Here we will take the approach of the method of separation of variables and assume a
solution of the form
H(x, y) = A(x) + B(y),
where A and B are functions to be determined. With this assumption, we get
dA k dB
y − x = 0.
dx m dy
Rearranging, we get
1 dA k dB
= .
x dx my dy
Now the term on the left is a function of x only, and the term on the right is a function of y only. The
only way this can be generally valid is if both terms are equal to the same constant, which we take to
be C. Hence,
1 dA k dB
= = C,
x dx my dy
from which we get two ordinary differential equations:
dA dB Cm
= Cx, = y.
dx dy k
The solution is
1 2 1 Cm 2
A(x) = Cx + K1 , B(y) = y + K2 .
2 2 k
A general solution is
1 2 m 2
H(x, y) = C x + y + K1 + K2 .
2 k
While this general solution is perfectly valid, we can obtain a common physical interpretation by taking
C = k, K1 + K2 = 0. With these choices, the Hamiltonian becomes
1 2 1
H(x, y) = kx + my 2 .
2 2
The first term represents the potential energy of the spring, the second term represents the kinetic
energy. Since by definition dH
dt = 0, this system conserves its mechanical energy. Verifying the properties
of a Hamiltonian, we see
dH ∂H dx ∂H dy
= + ,
dt ∂x dt ∂y dt
k
= kxy + my − x ,
m
= 0.
Since this system has dHdt = 0, then H(x, y) must be constant for all time, including t = 0, when the
initial conditions apply. So
1
H(x(t), y(t)) = H(x(0), y(0)) = kx20 + my02 .
2
Thus the system has the integral
1 1
kx2 + my 2 = kx20 + my02 .
2 2
Example 9.14
Using the Poincaré sphere, find the global dynamics, including at infinity, for the simple system
dx
= x, (9.65)
dt
dy
= −y. (9.66)
dt
Obviously the equilibrium point is at (x, y) = (0, 0), and that point is a saddle node.
Let us project the two state variables x and y into a three-dimensional space by the mapping
R2 → R3 :
x
X = p , (9.67)
1 + x2 + y 2
y
Y = p , (9.68)
1 + x2 + y 2
1
Z = p . (9.69)
1 + x2 + y 2
Note that
Note further if both x and y go to infinity, say on the line y = mx, then
1
lim X = √ , (9.72)
x→∞,y=mx m2 + 1
m
lim Y = √ , (9.73)
x→∞,y=mx m2 + 1
lim X2 + Y 2 = 1. (9.74)
x→∞,y=mx
4
Henri Poincaré, 1854-1912, French polymath.
So points at infinity are mapping onto a unit circle in (X, Y ) space. Also, going into the saddle node
at (x, y) = (0, 0) along the same line gives
So the original and transformed space have the same essential behavior near the finite equilibrium point.
Last, note that
x2 + y 2 + 1
X2 + Y 2 + Z2 = = 1. (9.77)
1 + x2 + y 2
Thus, in fact, the mapping takes one onto a unit sphere in (X, Y, Z) space. The surface X 2 +Y 2 +Z 2 = 1
is called the Poincaré sphere. One can actually view this in the same way one does an actual map of
the surface of the Earth. Just as a Mercator5 projection map is a representation of the spherical surface
of the earth projected onto a flat surface (and vice versa), the original (x, y) phase space is a planar
representation of the surface of the Poincaré sphere.
Let us find the inverse transformation. By inspection, it is seen that
X
x = , (9.78)
Z
Y
y = . (9.79)
Z
Now apply the transformation, Eqs. (9.78,9.79) to our dynamical system, Eqs. (9.65,9.66):
d X X
= , (9.80)
dt Z Z
d Y Y
= − . (9.81)
dt Z Z
Expand using the quotient rule to get
1 dX X dZ X
− 2 = , (9.82)
Z dt Z dt Z
1 dY Y dZ Y
− 2 = − . (9.83)
Z dt Z dt Z
Now on the unit sphere X 2 + Y 2 + Z 2 = 1, we must have
Y
1
0.75
0.5
0.25
X
-1 -0.75 -0.5 -0.25 0.25 0.5 0.75 1
-0.25
-0.5
-0.75
-1
Figure 9.8: Global phase portrait projection of the system dx/dt = x, dy/dt = −y on the
Poincaré sphere.
Example 9.15
Using projective space, find the global dynamics, including at infinity, for the same simple system
dx
= x, (9.106)
dt
dy
= −y. (9.107)
dt
Again, it is obvious that the equilibrium point is at (x, y) = (0, 0), and that point is a saddle node.
Let us project the two state variables x and y into a new two-dimensional space by the mapping
R2 → R2 :
1
X = , (9.108)
x
y
Y = . (9.109)
x
1
Y
−1
−2
−2 −1 0 1 2
X
Figure 9.9: Global phase portrait projection of the system dx/dt = x, dy/dt = −y on and
beyond the Poincaré sphere.
1
x = , (9.110)
X
Y
y = . (9.111)
X
Expanding, we find
1 dX 1
− = , (9.114)
X 2 dt X
1 dY Y dX Y
− 2 = − . (9.115)
X dt X dt X
Simplifying gives
dX
= −X, (9.116)
dt
dY dX
X −Y = −XY. (9.117)
dt dt
Solving for the derivatives, the system reduces to
dX
= −X, (9.118)
dt
dY
= −2Y. (9.119)
dt
By inspection, there is a sink at (X, Y ) = (0, 0). At such a point, the inverse mapping tells us x → ±∞
depending on whether X is positive or negative, and y is indeterminite. If we approach (X, Y ) = (0, 0)
along the line Y = mX, then y approaches the finite number m. This is consistent with trajectories
being swept away from the origin towards x → ±∞ in the original phase space, indicating an attraction
at x → ±∞. But it does not account for the trajectories emanating from y → ±∞. This is because
the transformation selected obscured this root.
To recover it, we can consider the alternate transformation X̂ = x/y, Ŷ = 1/y. Doing so leads to
the system dX̂/dt = 2X̂, dŶ /dt = Ŷ , which has a source at (X̂, Ŷ ) = (0, 0), which is consistent with
the source-like behavior in the original x, y space as y → ±∞. This transformation, however, obscures
the sink like behavior at x → ±∞.
To capture both points at infinity, we can consider a non-degenerate transformation, of which there
are infinitely many. One is X̃ = 1/(x + y), Ỹ = (x − y)/(x + y). Doing so leads to the system
dX̃/dt = −X̃ Ỹ , dỸ /dt = 1 − Ỹ 2 . This system has two roots, a source at (X̃, Ỹ ) = (0, −1) and a sink
at (X̃, Ỹ ) = (0, 1). The source corresponds to y → ±∞. The sink corresponds to x → ±∞.
9.8 Fractals
In the discussion on attractors in Section 9.6.1, we included geometrical shapes called frac-
tals. These are objects that are not smooth, but occur frequently in the dynamical systems
literature either as attractors or as boundaries of basins of attractions.
A fractal can be defined as a geometrical shape in which the parts are in some way similar
to the whole. This self-similarity may be exact, i.e. a piece of the fractal, if magnified, may
look exactly like the whole fractal. Before discussing examples we need to put forward
a working definition of dimension. Though there are many definitions in current use, we
present here the Hausdorff-Besicovitch6 dimension D. If Nǫ is the number of ‘boxes’ of side
length ǫ needed to cover an object, then
ln Nǫ
D = lim . (9.120)
ǫ→0 ln(1/ǫ)
We can check that this definition corresponds to the common geometrical shapes.
6
after Felix Hausdorff, 1868-1942, German mathematician, and Abram Samoilovitch Besicovitch, 1991-
1970, Russian mathematician.
A fractal has a dimension that is not an integer. Many physical objects are fractal-like, in
that they are fractal within a range of length scales. Coastlines are among the geographical
features that are of this shape. If there are Nǫ units of a measuring stick of length ǫ, the
measured length of the coastline will be of the power-law form ǫNǫ = ǫ1−D , where D is the
dimension.
k=0
k=1
k=2
k=3
k=4
ln Nǫ ln 2k k ln 2 ln 2
D = lim = lim k
= = = 0.6309 . . . . (9.121)
ǫ→0 ln(1/ǫ) k→∞ ln 3 k ln 3 ln 3
It can be seen that the endpoints of the removed intervals are never removed; it can be
shown the Cantor set contains an infinite number of points, and it is an uncountable set. It
is totally disconnected and has a Lebesgue measure zero.
and in the limit gives a continuous, closed curve that is nowhere smooth. Since Nǫ = 3 × 4k
and ǫ = 1/3k , the dimension of the Koch8 curve is
ln Nǫ ln(3)4k ln 3 + k ln 4 ln 4
D = lim = lim k
= lim = = 1.261 . . . . (9.122)
ǫ→0 ln(1/ǫ) k→∞ ln 3 k→∞ k ln 3 ln 3
The limit curve itself has infinite length, it is nowhere differentiable, and it surrounds a finite
area.
8
Niels Fabian Helge von Koch, 1870-1924, Swedish mathematician.
9
Karl Menger, 1902-1985, Austrian-born mathematician and active member of the influential
“Vienna Circle.” He served on the faculties of the Universities of Amsterdam, Vienna, Notre Dame, and the
Illinois Institute of Technology.
where a is real, b is odd, and ab > 1 + 3π/2. It is everywhere continuous, but nowhere
differentiable! Both require some effort to prove. A Weierstrass function is plotted in Figure
9.13. Its fractal character can be seen when one recognizes that cosine waves of ever higher
0.75
0.5
0.25
t
0.1 0.2 0.3 0.4 0.5
-0.25
-0.5
-0.75
stays bounded as k → ∞, when z0 = 0. The boundaries of this set are fractal. A Mandelbrot
set is sketched in Figure 9.14.
10
Karl Theodor Wilhalm Weierstrass, 1815-1897, Westphalia-born German mathematician.
11
Benoı̂t Mandelbrot, 1924-2010, Polish-born mathematician based mainly in France.
Figure 9.14: Mandelbrot set. Black regions stay bounded; colored regions become unbounded
with shade indicating how rapidly the system becomes unbounded. Image generated from
http://aleph0.clarku.edu/∼djoyce/cgi-bin/expl.cgi.
Associated with each c for the Mandelbrot set is a Julia12 set. In this case, the Julia set
is the set of complex initial seeds z0 which allow zk+1 = zk2 + c to converge for fixed complex
c. A Julia set for c = 0.49 + 0.57i is plotted in Figure 9.15.
9.9 Bifurcations
Dynamical systems representing some physical problem frequently have parameters associ-
ated with them. Thus, for x ∈ Rn , t ∈ R1 , λ ∈ R1 , f : Rn → Rn , we can write
dxi
= fi (x1 , x2 , · · · , xn ; λ) (i = 1, · · · , n), (9.125)
dt
where λ is a parameter. The theory can easily be extended if there is more than one
parameter.
We would like to consider the changes in the behavior of t → ∞ solutions as the real
number λ, called the bifurcation parameter, is varied. The nature of the critical point may
change as the parameter λ is varied; other critical points may appear or disappear, or its
stability may change. This is a bifurcation, and the λ at which it happens is the bifurcation
point. The study of the solutions and bifurcations of the steady state falls under singularity
theory.
12
Gaston Maurice Julia, 1893-1978, Algerian-born French mathematician.
Figure 9.15: Julia set for c = 0.49+0.57i. Black regions stay bounded; colored regions become
unbounded with shade of color indicating how rapidly the system becomes unbounded. Image
generated from http://aleph0.clarku.edu/∼djoyce/cgi-bin/expl.cgi.
Let us look at some of the bifurcations obtained for different vector fields. Some of the
examples will be one-dimensional, i.e. x ∈ R1 , λ ∈ R1 , f : R1 → R1 .
dx
= f (x; λ). (9.126)
dt
Even though this can be solved exactly in most cases, we will assume that such a solution
is not available so that the techniques of analysis can be developed for more complicated
systems. For a coefficient matrix that is a scalar, the eigenvalue is the coefficient itself. The
eigenvalue will be real and will cross the imaginary axis of the complex plane through the
origin as λ is changed. This is called a simple bifurcation.
For λ < λ0 , the critical point is asymptotically stable; for λ > λ0 it is unstable.
Notice that the function V (x) = x2 satisfies the following conditions: V > 0 for x 6= 0,
V = 0 for x = 0, and dV dt
= dV dx
dx dt
= −2x2 (x2 − (λ − λ0 )) ≤ 0 for λ < λ0 . Thus V (x) is a
Lyapunov function and x = 0 is globally stable for all perturbations, large or small, as long
as λ < λ0 . √
Now let us examine the critical point x = λ − λ0 which exists only for λ > λ0 . Putting
x = x + x′ , the right side of equation (9.127) becomes
p p 2
′ ′
f (x) = − λ − λ0 + x λ − λ0 + x − (λ − λ0 ) .
λο λ
dx
= −x3 . (9.128)
dt
CC BY-NC-ND. 28 March 2011, M. Sen, J. M. Powers.
9.9. BIFURCATIONS 397
This equation has a critical point at x = 0 but has no linearization. We must do a non-linear
analysis to determine the stability of the critical point. In this case it is straightforward.
Solving directly and applying an initial condition, we obtain
x(0)
x(t) = ± p , (9.129)
1 + 2x(0)2 t
lim x(t) = 0. (9.130)
t→∞
Since the system approaches the critical point as t → ∞ for all values of x(0), the critical
point x = 0 unconditionally stable.
λο λ
λο λ
Example 9.16
With x, y, t, λ, λ0 ∈ R1 , take
dx
= (λ − λ0 )x − y − x(x2 + y 2 ),
dt
dy
= x + (λ − λ0 )y − y(x2 + y 2 ).
dt
The origin (0,0) is a critical point. The linearized perturbation equations are
′ ′
d x λ − λ0 −1 x
= .
dt y′ 1 λ − λ0 y′
The eigenvalues µ of the coefficient matrix are µ = (λ − λ0 ) ± i. For λ < λ0 the real part is negative
and the origin is stable. At λ = λ0 there is a Hopf13 bifurcation as the eigenvalues cross the imaginary
axis of the complex plane as λ is changed. For λ > λ0 a periodic orbit in the (x, y) phase plane appears.
The linear analysis will not give the amplitude of the motion. Writing the given equation in polar
coordinates (r, θ)
dr
= r(λ − λ0 ) − r3 ,
dt
dθ
= 1.
dt
This is a pitchfork bifurcation in the amplitude of the oscillation r.
13
Eberhard Frederich Ferdinand Hopf, 1902-1983, Austrian-born, German mathematician.
y − x = 0, (9.144)
rx − y − xz = 0, (9.145)
−bz + xy = 0, (9.146)
which give p p
x 0 p b(r − 1) −pb(r − 1)
y = 0 , b(r − 1) , − b(r − 1) . (9.147)
z 0 r−1 r−1
A linear stability analysis of each critical point follows.
Thus the real parts of the eigenvalues are negative if r < rc = σ(σ+b+3)
σ−b−1
. At r = rc the
characteristic equation (9.152) can be factored to give the eigenvalues −(σ + b+ 1), and
±i 2σ(σ+1)
σ−b−1
, corresponding to a Hopf bifurcation. The periodic solution which is created
at this value of r can be shown to be unstable so that the bifurcation is subcritical.
Example 9.17
Consider the solution to the Lorenz equations for conditions: σ = 1, r = 28, b = 8/3 with initial
conditions x(0) = y(0) = z(0) = 1. The fixed point is given by
r
p 8
x = b(r − 1) = (28 − 1) = 8.485, (9.158)
3
r
p 8
y = b(r − 1) = (28 − 1) = 8.485, (9.159)
3
z = r − 1 = 28 − 1 = 27. (9.160)
Consideration of the roots of the characteristic equation shows the fixed point here is stable:
14 2 232
λ3 + λ + λ + 144 = 0, (9.162)
3 3
√
4 2528
λ = −2, λ=− ± i. (9.163)
3 6
Figure 9.19 shows the phase space trajectories in x, y, z space and the behavior in the time domain,
x(t), y(t), z(t).
x
10
8
x 6
4 8 4
20
y 2
10
t
0 0 1 2 3 4 5 6
40 y
30
30 20
z 10
20
t
1 2 3 4 5 6
10
z
50
40
30
20
10
t
0 1 2 3 4 5 6
Figure 9.19: Solution to Lorenz equations, σ = 1, r = 28, b = 8/3. Initial conditions are
x(0) = y(0) = z(0) = 1.
Example 9.18
Now consider the conditions: σ = 10, r = 28, b = 8/3. Initial conditions are x(0) = y(0) = z(0) = 1.
The fixed point is again given by
r
p 8
x = b(r − 1) = (28 − 1) = 8.485, (9.164)
3
r
p 8
y = b(r − 1) = (28 − 1) = 8.485, (9.165)
3
z = r − 1 = 28 − 1 = 27. (9.166)
Now, consideration of the roots of the characteristic equation shows the fixed point here is unstable:
x
20
10
20 0 t
y 5 10 15 20 25
0 -10
-20
y
40 20
30
z 0 t
20 5 10 15 20 25
10 -20
0
z
-10
0 40
x 10 20
20
0 t
5 10 15 20 25
Figure 9.20: Phase space trajectory and time domain plots for solution to Lorenz equations,
σ = 10, r = 28, b = 8/3. Initial conditions are x(0) = y(0) = z(0) = 1.
du dv
+σ = −σ(1 + σ)v, (9.177)
dt dt
du dv
− = (1 + σ)v − (u + σv)w, (9.178)
dt dt
CC BY-NC-ND. 28 March 2011, M. Sen, J. M. Powers.
9.10. LORENZ EQUATIONS 405
dw
= −bw + (u + σv)(u − v). (9.179)
dt
Solving directly for the derivatives so as to place the equations in autonomous form, we get
du σ
= 0u − (u + σv)w = λ1 u + nonlinear terms, (9.180)
dt 1+σ
dv 1
= −(1 + σ)v + (u + σv)w = λ2 v + nonlinear terms, (9.181)
dt 1+σ
dw
= −bw + (u + σv)(u − v) = λ3 w + nonlinear terms. (9.182)
dt
The objective of using the eigenvectors as basis vectors is to change the original system to
diagonal form in the linear terms. Notice that the linear portion of the system is in diagonal
form with the coefficients on each linear term as a distinct eigenvalue. Furthermore, the
eigenvalues λ2 = −(1 + σ) and λ3 = −b are negative ensuring that the linear behavior
v = e−(1+σ)t and w = e−bt takes the solution very quickly to zero in these variables.
It would appear then that we are only left with an equation in u(t) for large t. However,
if we put v = w = 0 in the right side, dv/dt and dw/dt would be zero if it were not for the
u2 term in dw/dt, implying that the dynamics is confined to v = w = 0 only if we ignore
this term. According to the center manifold theorem it is possible to find a line (called the
center manifold) which is tangent to u = 0, but is not necessarily the tangent itself, to which
the dynamics is indeed confined.
We can get as good an approximation to the center manifold as we want by choosing new
variables. Expanding the equation for dw dt
, which has the potential problem, we get
dw
= −bw + u2 + (σ − 1)uv − σv 2 . (9.183)
dt
Letting
u2
w̃ = w − , (9.184)
b
so that −bw + u2 = −bw̃, we can eliminate the potential problem with the derivative of w.
In the new variables (u, v, w̃), the full Lorenz equations are written as
du σ u2
= − (u + σv)(w̃ + ), (9.185)
dt 1+σ b
dv 1 u2
= −(1 + σ)v + (u + σv)(w̃ + ), (9.186)
dt 1+σ b
dw̃ 2σ u2
= −bw̃ + (σ − 1)uv − σv 2 + u(u + σv)(w̃ + ). (9.187)
dt b(1 + σ) b
Once again the variables v and w̃ go to zero very quickly. Formally setting them to zero,
and examining all equations we see that
du σ
= − u3 , (9.188)
dt b(1 + σ)
dv 1
= u3 , (9.189)
dt (1 + σ)b
dw̃ 2σ
= u4 . (9.190)
dt (1 + σ)b2
Here dv/dt and dw̃/dt approach zero if u approaches zero. Now the equation for the evolution
of u suggests that this is the case. Simply integrating this equation and applying and initial
condition we get
s
b(1 + σ)
u(t) = ±(u(0)) (9.191)
b(1 + σ) + 2σ(u(0))2t
center
0.3 manifold
0.2
0.1
.
stable
equilibrium
0.2 0.4 0.6 0.8 1
u
point
Figure 9.21: Projection of solution trajectory and center manifold for forced Lorenz equations
at bifurcation point.
Problems
1. For the logistics equation: xk+1 = rxk (1 − xk ); 0 < xk < 1, 0 < r < 4, write a short program which
determines the value of x as k → ∞. Plot the bifurcation diagram, that is the limiting value of x as
a function of r for 0 < r < 4. If ri is the ith bifurcation point, that is the value at which the number
of fixed points changes, make an estimate of Feigenbaum’s constant,
ri−1 − ri
δ = lim
i→∞ ri − ri+1
2. If
dx dy
x + xy =x−1
dt dt
dx dy
(x + y) +x =y+1
dt dt
write the equation in autonomous form,
dx
= f (x, y)
dt
dy
= g(x, y)
dt
Plot the lines f = 0, g = 0 in the xy phase plane. Also plot in this plane the vector field defined by
the differential equations. With a combination of analysis and numerics, find a path in phase space
from one critical point to the other. For this path, plot x(t), y(t) and include the path in the xy phase
plane.
dx
= x3 + x (λ − 3)2 − 1
dt
where λ is the bifurcation parameter, indicating stable and unstable branches.
5. A two-dimensional dynamical system expressed in polar form is
dr
= r(r − 2)(r − 3)
dt
dθ
= 2
dt
Find the (a) critical point(s), (b) periodic solution(s), and (c) analyze their stability.
6. Find a critical point of the following system, and show its local and global stability.
dx
= (x − 2) (y − 1)2 − 1
dt
dy
= (2 − y) (x − 2)2 + 1
dt
dz
= (4 − z)
dt
7. Find the general solution of dx/dt = A · x where
1 −3 1
A = 2 −1 −2
2 −3 0
18. Let
1 1 2
A= 0 1 1
0 0 1
Solve the equation
dx
= A · x.
dt
Determine the critical points and their stability.
19. Draw the bifurcation diagram of
dx
= (x2 − 2)2 − 2(x2 + 1)(λ − 1) + (λ − 1)2
dt
indicating the stability of each branch.
20. Show that for all initial conditions the solutions of
dx
= −x + x2 y − y 2
dt
dy
= −x3 + xy − 6z
dt
dz
= 2y
dt
tend to x = y = z = 0 as t → ∞.
21. Draw the bifurcation diagram of
dx
= x3 + x (λ − 2)3 − 1
dt
where λ is the bifurcation parameter, indicating stable and unstable branches.
22. Solve the system of equations dx/dt = A · x where
−3 0 2 0
0 −2 0 0
A= 0
0 1 1
0 0 0 0
24. Analyze the local stability of the origin in the following system
dx
= −2x + y + 3z + 8y 3
dt
dy
= −6y − 5z + 2z 3
dt
dz
= z + x2 + y 3 .
dt
27. Find the dynamical system corresponding to the Hamiltonian H(x, y) = x2 + 2xy + y 2 and then solve
it.
28. Show that solutions of the system of differential equations
dx
= −x + y 3 − z 3
dt
dy
= = −y + z 3 − x3
dt
dz
= −z + x3 − y 3
dt
eventually approach the origin for all initial conditions.
29. Find and sketch all critical points (x, y) of
dx
= (λ − 1)x − 3xy 2 − x3
dt
dy
= (λ − 1)y − 3x2 y − y 3
dt
as functions of λ. Determine the stability of (x, y) = (0, 0), and of one post-bifurcation branch.
30. Write in matrix form and solve
dx
= y+z
dt
dy
= z+x
dt
dz
= x+y
dt
31. Find the critical point (or points) of the Van der Pol equation
d2 x dx
2
− a(1 − x2 ) + x = 0, a > 0
dt dt
and determine its (or their) stability to small perturbations. For a = 1, plot the dx/dt, x phase plane
including critical points and vector fields.
32. Consider a straight line between x = 0 and x = l. Remove the middle half (i.e. the portion between
x = l/4 and x = 3l/4). Repeat the process on the two pieces that are left. Find the dimension of
what is left after an infinite number of iterations.
35. Find a Lyapunov function of the form V = ax2 + by 2 to investigate the global stability of the critical
point x = y = 0 of the system of equations
dx
= −2x3 + 3xy 2
dt
dy
= −x2 y − y 3
dt
36. Draw a bifurcation diagram for the differential equation
dx
= (x − 3)(x2 − λ)
dt
Analyze linear stability and indicate stable and unstable branches.
37. Solve the following system of differential equations using generalized eigenvectors
dx
= −5x + 2y + z
dt
dy
= −5y + 3z
dt
dz
= = −5z
dt
38. Analyze the linear stability of the critical point of
dx
= 2y + y 2
dt
dy
= −λ + 2x2
dt
39. Show that the solutions of
dx
= y − x3
dt
dy
= −x − y 3
dt
tend to (0,0) as t → ∞.
40. Sketch the bifurcation diagram showing the stable and unstable steady states of
dx
= λx(1 − x) − x
dt
dx
= a + x2 y − 2bx − x
dt
dy
= bx − x2 y
dt
where a, b > 0.
42. Show that the Hénon-Heiles system
d2 x
= −x − 2xy
dt2
d2 y
= −y + y 2 − x2
dt2
is Hamiltonian. Find the Hamiltonian of the system, and determine the stability of the critical point
at the origin.
43. Solve dx/dt = A · x where
2 1
A=
0 2
using the exponential matrix.
44. Sketch the steady state bifurcation diagrams of
dx
= (x − λ)(x + λ)((x − 3)2 + (λ − 1)2 − 1)
dt
Determine the linear stability of each branch; indicate the stable and unstable ones differently on the
diagram.
45. Classify the critical point of
d2 x
+ (λ − λ0 )x = 0
dt2
46. Show that x = 0 is a stable critical point of the differential equation
dx X
=− nai x2i+1
dt i=0
where ai ≥ 0, i = 0, 1, · · · , n.
47. Find the stability of the critical points of the Duffing equation
d2 x dx
=a − bx + x3 = 0
dt2 dt
for positive and negative values of a and b. Sketch the flow lines.
48. Find a Lyapunov function to investigate the critical point x = y = 0 of the system of equations
dx
= −2x3 + 3xy 2
dt
dx
= −x2 y − y 3
dt
49. The populations x and y of two competing animal species are governed by
dx
= x − 2xy
dt
dy
= −y + xy
dt
What are the steady-state populations? Is the situation stable?
50. For the Lorenz equations with b = 8/3, r = 28 and initial conditions x(0) = 2, y(0) = 1, z(0) = 3,
numerically integrate the Lorenz equations for two cases, σ = 1, σ = 10. For each case plot the
trajectory in xyz phase space and plot x(t), y(t), z(t) for 0 < t < 50. Change the initial condition on
x to x(0) = 2.002 and plot the difference in the predictions of x versus time for both values of σ.
51. Use the Poincaré sphere to find all critical points, finite and infinite of the system
dx
= 2x − 2xy
dt
dy
= 2y − x2 + y 2
dt
Plot families of trajectories in the x, y phase space and the X, Y projection of the Poincaré sphere.
Appendix
The material in this section is not covered in detail; some is review from undergraduate
classes.
1 1
sin x sin y = cos(x − y) − cos(x + y)
2 2
1 1
sin x cos y = sin(x + y) + sin(x − y)
2 2
1 1
cos x cos y = cos(x − y) + cos(x + y)
2 2
1 1
sin2 x = − cos 2x
2 2
1
sin x cos x = sin 2x
2
1 1
cos2 x = + cos 2x
2 2
3 1
sin3 x = sin x − sin 3x
4 4
1 1
sin2 x cos x = cos x − cos 3x
4 4
1 1
sin x cos2 x = sin x + sin 3x
4 4
3 1
cos3 x = cos x + cos 3x
4 4
3 1 1
sin4 x = − cos 2x + cos 4x
8 2 8
415
416 CHAPTER 10. APPENDIX
1 1
sin3 x cos x = sin 2x − sin 4x
4 8
1 1
sin2 x cos2 x = − cos 4x
8 8
1 1
sin x cos3 x = sin 2x + sin 4x
4 8
3 1 1
cos4 x = + cos 2x + cos 4x
8 2 8
5 5 1
sin5 x = sin x − sin 3x + sin 5x
8 16 16
1 3 1
sin4 x cos x = cos x − cos 3x + cos 5x
8 16 16
1 1 1
sin3 x cos2 x = sin x + sin 3x − sin 5x
8 16 16
1 1 1
sin2 x cos3 x = − cos x − cos 3x − cos 5x
8 16 16
1 3 1
sin x cos4 x = sin x + sin 3x + sin 5x
8 16 16
5 5 1
cos5 x = cos x + cos 3x + cos 5x
8 16 16
a0 sn + a1 sn−1 + . . . + an−1 s + an = 0
has roots with negative real parts if and only if the following conditions are satisfied:
(i) a1 /a0 , a2 /a0 , . . . , an /a0 > 0
(ii) Di > 0, i = 1, . . . , n
The Hurwitz determinants Di are defined by
D1 = a1
a a3
D2 = 1
a0 a2
a1 a3 a5
D3 = a0 a2 a4
0 a1 a3
1
Edward John Routh, 1831-1907, Canadian-born English mathematician, and Adolf Hurwitz, 1859-
1919, German mathematician.
a1 a3 a5 ... a2n−1
a0 a2 a4 ... a2n−2
0 a1 a3 ... a2n−3
Dn = 0 a0 a2 ... a2n−4
.. .. .. .. ..
. . . . .
0 0 0 ... an
with ai = 0, if i > n.
1
f (x) = f (a) + f ′ (a)(x − a) + f ′′ (a)(x − a)2 + · · · (10.2)
2
where the function and its derivatives on the right side are evaluated at x = a. This is a
power series for f (x). We have used primes to indicate derivatives.
Example 10.1
Expand (1 + x)n about x = 0.
f (x) = (1 + x)n
f (0) = 1
f ′ (0) = n
f ′′ (0) = n(n − 1)
..
.
1
(1 + x)n = 1 + nx + n(n − 1)x2 + · · ·
2
Bender and Orszag show that Stirling’s2 formula is a divergent series. It is an asymptotic
series, but as more terms are added, the solution can actually get worse. The Gamma
function and its amplitude are plotted in Figure 10.1.
2
James Stirling, 1692-1770, Scottish mathematician and member of a prominent Jacobite family.
Γ
15 Γ
10 1. x 10 8
5 10000
x 1
-4 -2 2 4
-5 0.0001
-10 1. x 10 - 8
-15 x
-15 -10 -5 0 5 10 15
Property:
Γ(p)Γ(q)
B(p, q) = (10.6)
Γ(p + q)
ζ
ζ
10
8 100000.
6 100
4 0.1
2 0.0001
x x
2 4 6 8 10 -30 -20 -10 0 10
Figure 10.2: Riemann zeta function and amplitude of Riemann zeta function.
erf (x) erfc (x)
1 2
0.5 1.5
1
-4 -2 2 4 x
0.5
-0.5
-1 -4 -2 2 4 x
C(x) S(x)
0.75
0.6
0.5
0.4
0.25
0.2
-0.5 -0.4
-0.6
-0.75
The sine integral function is real valued for x ∈ (−∞, ∞). The cosine integral function is
real valued for x ∈ [0, ∞). We also have limx→0+ Ci(x) → −∞. The cosine integral takes on
a value of zero at discrete positive real values, and has an amplitude which slowly decays as
x → ∞. The sine integral and cosine integral functions are plotted in Figure 10.5.
Si (x) Ci (x)
1.5 0.4
1
0.2
0.5
-20 -10 10 20 x 5 10 15 20 25 30
-0.5 x
-1 -0.2
-1.5
Figure 10.5: Sine integral function, Si(x), and cosine integral function Ci(x).
Another common way of writing the elliptic integral is to take η = sin φ, so that
Z φ
dφ
F (φ, k) = p (10.15)
0 (1 − k 2 sin2 φ)
The Legendre elliptic integral of the second kind is
Z y
(1 − k 2 η 2 )
E(y, k) = p dη (10.16)
0 (1 − η 2 )
which is equivalent to Z φ q
Π(φ, n, k) = 1 − k 2 sin2 φ dφ (10.19)
0
For φ = π/2, we have the complete elliptic integrals:
π Z π/2
dφ
F ,k = p (10.20)
2 0 (1 − k 2 sin2 φ)
π Z π/2 q
E ,k = 1 − k 2 sin2 φ dφ (10.21)
2 0
π Z π/2 q
Π , n, k = 1 − k 2 sin2 φ dφ (10.22)
2 0
δ(x − a) = lim+ ∆ǫ (x − a)
ǫ→0
1.
0 if x < a − 2ǫ
1
∆ǫ (x − a) = if a − 2ǫ ≤ x ≤ a + ǫ
(10.27)
ǫ 2
0 if x > a + 2ǫ
2.
ǫ
∆ǫ (x − a) = (10.28)
π((x − a)2 + ǫ2 )
3.
1 2
∆ǫ (x − a) = √ e−(x−a) /ǫ (10.29)
πǫ
then
d
H(x − a) = δ(x − a) (10.32)
dx
The generator of the Dirac function ∆ǫ (x − a) and the generator of the Heaviside function
h(x−a) are plotted for a = 0 and ǫ = 1/5 in Figure 10.6. As ǫ → 0, ∆ǫ has its width decrease
and its height increase in such a fashion that its area remains constant; simultaneously h
has its slope steepen in the region where it jumps from zero to unity as ǫ → 0.
∂f ∂f ∂f
df = dx1 + dx2 + · · · + dxn . (10.33)
∂x1 ∂x2 ∂xn
5
Oliver Heaviside, 1850-1925, English mathematician.
∆ ε (x) h (x)
5 1
Dirac Heaviside
4 0.8
delta function step function
generator 3 generator 0.6
2 0.4
1 0.2
Figure 10.6: Generators of Dirac delta function and Heaviside function, ∆ǫ (x − a) and
h(x − a) plotted for a = 0 and ǫ = 1/5.
1 1 1 1
eiθ = 1 + iθ + (iθ)2 + (iθ)3 + (iθ)4 + (iθ)5 + . . . , (10.41)
2! 3! 4! 5!
1 2 1 3 1 4 1 5
= 1 + iθ − θ − i θ + θ + i θ + . . . (10.42)
2! 3! 4! 5!
As the two series are identical, we have Euler’s formula
z = x + iy, (10.45)
p
we can multiply and divide by x2 + y 2 to obtain
!
p x y
z= x2 + y 2 p + ip . (10.46)
2
x +y 2 x + y2
2
Noting the similarities between this and the transformation between Cartesian and polar
coordinates suggests we adopt
p x y
r = x2 + y 2 , cos θ = p , sin θ = p . (10.47)
x2 + y 2 x2 + y 2
Thus we have
The polar and Cartesian representation of a complex number z is shown in Figure 10.7.
Now we can define the complex conjugate z as
z = x − iy, (10.50)
!
p x y
z = x2 + y 2 p − ip , (10.51)
x2 + y 2 x2 + y 2
z = r (cos θ − i sin θ) , (10.52)
z = r (cos(−θ) + i sin(−θ)) , (10.53)
z = re−iθ . (10.54)
6
Abraham de Moivre, 1667-1754, French mathematician.
iy
y
2
2 +y
x
r=
x x
Now there are many paths that we can choose to evaluate the derivative. Let us consider
two distinct paths, y = C1 and x = C2 . We will get a result which can be shown to be valid
for arbitrary paths. For y = C1 , we have ∆z = ∆x, so
dW W (xo + iyo + ∆x) − W (xo + iyo ) ∂W
= = . (10.65)
dz zo ∆x ∂x y
For x = C2 , we have ∆z = i∆y, so
dW W (xo + iyo + i∆y) − W (xo + iyo ) 1 ∂W ∂W
= = = −i . (10.66)
dz zo i∆y i ∂y x ∂y x
Now for an analytic function, we need
∂W ∂W
= −i . (10.67)
∂x y ∂y x
or, expanding, we need
∂φ ∂ψ ∂φ ∂ψ
+i = −i +i , (10.68)
∂x ∂x ∂y ∂y
∂ψ ∂φ
= −i . (10.69)
∂y ∂y
For equality, and thus path independence of the derivative, we require
∂φ ∂ψ ∂φ ∂ψ
= , =− . (10.70)
∂x ∂y ∂y ∂x
These are the well known Cauchy-Riemann equations for analytic functions of complex
variables.
Now most common functions are easily shown to be analytic. For example for the function
W (z) = z 2 + z, which can be expressed as W (z) = (x2 + x − y 2) + i(2xy + y), we have
φ(x, y) = x2 + x − y 2 , ψ(x, y) = 2xy + y, (10.71)
∂φ ∂ψ
= 2x + 1, = 2y, (10.72)
∂x ∂x
∂φ ∂ψ
= −2y, = 2x + 1. (10.73)
∂y ∂y
Note that the Cauchy-Riemann equations are satisfied since ∂φ∂x
= ∂ψ
∂y
and ∂φ
∂y
= − ∂ψ
∂x
. So the
derivative is independent of direction, and we can say
dW ∂W
= = (2x + 1) + i(2y) = 2(x + iy) + 1 = 2z + 1. (10.74)
dz ∂x y
We could get this result by ordinary rules of derivatives for real functions.
For an example of a non-analytic function consider W (z) = z. Thus
W (z) = x − iy. (10.75)
∂φ ∂φ ∂ψ ∂ψ ∂φ ∂ψ
So φ = x and ψ = −y, = 1,
∂x
= 0, and
∂y
= 0,
∂x
= −1. Since
∂y
6=∂x ∂y
, the
Cauchy-Riemann equations are not satisfied, and the derivative depends on direction.
Problems
1. Find the limit as x → 0 of
4 cosh x + sinh(arctan ln cos 2x) − 4
√ .
e−x + arcsin x − 1 + x2
dφ
2. Find dx in two different ways, where
Z x4
√
φ= x ydy.
x2
3. Determine
√
(a) 4 i
√
(b) ii i i
4. Write three terms of a Taylor series expansion for the function f (x) = exp(tan x) about the point
x = π/4. For what range of x is the series convergent?
5. Find all complex numbers z = x + iy such that |z + 2i| = |1 + i|.
3
6. Determine limn→∞ zn for zn = n + ((n + 1)/(n + 2))i.
7. A particle is constrained to a path which is defined by the function s(x, y) = x2 + y − 5 = 0. The
velocity component in the x-direction, dx/dt = 2y. What are the position and velocity components
in the y-direction when x = 4.
Rx 2
8. The error function is defined as erf (x) = √2π 0 e−u du. Determine its derivative with respect to x.
9. Verify that Z 2π
sin nx
lim dx = 0.
n→∞ π nx
10. Write a Taylor series expansion for the function f (x, y) = x2 cos y about the point x = 2, y = π.
Include the x2 , y 2 and xy terms.
11. Show that Z ∞
2
φ= e−t cos 2tx dt
0
satisfies the differential equation
dφ
+ 2φx = 0.
dx
12. Evaluate the Dirichlet discontinuous integral
Z
1 ∞ sin ax
I= dx
π −∞ x
for −∞ < a < ∞. You can use the results of example 3.11, Greenberg.
13. Defining
x3 − y 3
u(x, y) = ,
x2 + y 2
except at x = y = 0, where u = 0, show that ux (x, y) exists at x = y = 0 but is not continuous there.
14. Using complex numbers show that
(a) cos3 x = 41 (cos 3x + 3 cos x)
(b) sin3 x = 41 (3 sin x − sin 3x)
M. Abramowitz and I. A. Stegun, eds., Handbook of Mathematical Functions, Dover, New York, 1964.
V. I. Arnold, Ordinary Differential Equations, MIT Press, Cambridge, MA, 1973.
V. I. Arnold, Geometrical Methods in the Theory of Ordinary Differential Equations, Springer, New York,
NY, 1983.
A. A. Andronov, Qualitative Theory of Second Order Dynamical Systems, Wiley, New York, NY, 1973.
R. Aris, Vectors, Tensors, and the Basic Equations of Fluid Mechanics, Dover, New York, NY, 1962.
N. H. Asmar, Applied Complex Analysis with Partial Differential Equations, Prentice-Hall, Upper Saddle
River, NJ, 2002.
G. I. Barenblatt, Scaling, Self-Similarity, and Intermediate Asymptotics, Cambridge University Press, Cam-
bridge, UK, 1996.
R. Bellman and K. L. Cooke, Differential-Difference Equations, Academic Press, New York, NY, 1963.
C. M. Bender and S. A. Orszag, Advanced Mathematical Methods for Scientists and Engineers, Springer-
Verlag, New York, NY, 1999.
M. L. Boas, Mathematical Methods in the Physical Sciences, Third Edition, Wiley, New York, NY, 2005.
A. I. Borisenko and I. E. Tarapov, Vector and Tensor Analysis with Applications, Dover, New York, NY,
1968.
M. Braun, Differential Equations and Their Applications, Springer-Verlag, New York, NY, 1983.
I. N. Bronshtein and K. A. Semendyayev, Handbook of Mathematics, Springer, Berlin, 1998.
C. Canuto, M. Y. Hussaini, A. Quarteroni, and T. A. Zang, Spectral Methods in Fluid Dynamics, Springer-
Verlag, New York, NY, 1988.
G. F. Carrier and C. E. Pearson, Ordinary Differential Equations, SIAM, Philadelphia, PA, 1991.
P. G. Ciarlet, Introduction to Numerical Linear Algebra and Optimisation, Cambridge University Press,
Cambridge, UK, 1989.
J. A. Cochran, H. C. Wiser and B. J. Rice, Advanced Engineering Mathematics, Second Edition, Brooks/Cole,
Monterey, CA, 1987.
R. Courant and D. Hilbert, Methods of Mathematical Physics, Vols. 1 and 2, Wiley, New York, NY, 1989.
I. Daubechies, Ten Lectures on Wavelets, SIAM, Philadelphia, PA, 1992.
L. Debnath and P. Mikusinski, Introduction to Hilbert Spaces with Applications, Third Edition, Elsevier,
Amsterdam, Netherlands, 2005.
P. G. Drazin, Nonlinear Systems, Cambridge University Press, Cambridge, UK, 1992.
429
R. D. Driver, Ordinary and Delay Differential Equations, Springer-Verlag, New York, NY, 1977.
J. Feder, Fractals, Plenum Press, New York, NY, 1988.
B. A. Finlayson, The Method of Weighted Residuals and Variational Principles, Academic Press, New York,
NY, 1972.
B. Fornberg, A Practical Guide to Pseudospectral Methods, Cambridge, New York, NY, 1998.
B. Friedman, Principles and Techniques of Applied Mathematics, Dover Publications, New York, NY, 1956.
I. M. Gelfand and S. V. Fomin, Calculus of Variations, Dover, New York, NY, 2000.
J. Gleick, Chaos, Viking, New York, NY, 1987.
G. H. Golub and C. F. Van Loan, Matrix Computations, Third Edition, Johns Hopkins, Baltimore, MD,
1996.
S. W. Goode, An Introduction to Differential Equations and Linear Algebra, Prentice-Hall, Englewood Cliffs,
NJ, 1991.
D. Gottlieb and S. A. Orszag, Numerical Analysis of Spectral Methods: Theory and Applications, SIAM,
Philadelphia, PA, 1977.
M. D. Greenberg, Foundations of Applied Mathematics, Prentice-Hall, Englewood Cliffs, NJ, 1978.
J. Guckenheimer and P. H. Holmes, Nonlinear Oscillations, Dynamical Systems, and Bifurcations of Vector
Fields, Springer-Verlag, New York, NY, 1983.
J. Hale and H. Koçak, Dynamics and Bifurcations, Springer-Verlag, New York, NY, 1991.
F. B. Hildebrand, Advanced Calculus for Applications, 2nd Ed., Prentice-Hall, Englewood Cliffs, NJ, 1976.
M. W. Hirsch and S. Smale, Differential Equations, Dynamical Systems, and Linear Algebra, Academic
Press, Boston, MA, 1974.
M. H. Holmes, Introduction to Perturbation Methods, Springer-Verlag, New York, NY, 1995.
M. H. Holmes, Introduction to the Foundations of Applied Mathematics, Springer-Verlag, New York, NY,
2009.
M. Humi and W. Miller, Second Course in Ordinary Differential Equations for Scientists and Engineers,
Springer-Verlag, New York, NY, 1988.
E. J. Hinch, Perturbation Methods, Cambridge, Cambridge, UK, 1991.
D. W. Jordan and P. Smith, Nonlinear Ordinary Differential Equations, Clarendon Press, Oxford, UK, 1987.
P. B. Kahn, Mathematical Methods for Engineers and Scientists, Dover, New York, NY, 2004.
W. Kaplan, Advanced Calculus, Fifth Edition, Addison-Wesley, Boston, MA, 2003.
J. Kevorkian and J. D. Cole, Perturbation Methods in Applied Mathematics, Springer-Verlag, New York,
NY, 1981.
J. Kevorkian and J. D. Cole, Multiple Scale and Singular Perturbation Methods, Springer-Verlag, New York,
NY, 1996.
A. N. Kolmogorov and S. V. Fomin, Elements of the Theory of Functions and Functional Analysis, Dover,
New York, NY, 1999.
L. D. Kovach, Advanced Engineering Mathematics, Addison-Wesley, Reading, MA,1982.
430
E. Kreyszig, Advanced Engineering Mathematics, Ninth Edition, Wiley, New York, NY, 2005.
E. Kreyszig, Introductory Functional Analysis with Applications, Wiley, New York, NY, 1978.
P. D. Lax, Functional Analysis, Wiley, New York, NY, 2002.
P. D. Lax, Linear Algebra and its Applications, Second Edition, Wiley, Hoboken, NJ, 2007.
A. J. Lichtenberg and M. A. Lieberman, Regular and Chaotic Dynamics, Second Edition, Springer, Berlin,
1992.
C. C. Lin and L. A. Segel, Mathematics Applied to Deterministic Problems in the Natural Sciences, SIAM,
Philadelphia, PA, 1988.
J. R. Lee, Advanced Calculus with Linear Analysis, Academic Press, New York, NY, 1972.
J. D. Logan, Applied Mathematics, Third Edition, Wiley, Hoboken, NJ, 2006.
R. J. Lopez, Advanced Engineering Mathematics, Addison Wesley Longman, Boston, MA, 2001.
J. Mathews and R. L. Walker, Mathematical Methods of Physics, Addison-Wesley, Redwood City, CA, 1970.
431
S. Wiggins, Introduction to Applied Nonlinear Dynamical Systems and Chaos, Springer Verlag, New York,
NY, 1990.
M. Van Dyke, Perturbation Methods in Fluid Mechanics, Parabolic Press, Stanford, CA, 1975.
C. R. Wylie and L. C. Barrett, Advanced Engineering Mathematics, 6th Ed., McGraw-Hill, New York, NY,
1995.
E. Zeidler, Applied Functional Analysis, Springer Verlag, New York, NY, 1995.
D. G. Zill and M. R. Cullen, Advanced Engineering Mathematics, Third Edition, Jones and Bartlett, Boston,
MA, 2006.
432